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multiple integrals, etc.). I hope that concluding my book in this manner will encourage
the readers research on multivariable calculus, even if he or she does not plan to take the
course in college.
Before reading the book
I have generally listed the material with which the student should be familiar
before he or she begins reading the book:
I.) Limits and Continuity: This includes understanding what a limit is and how to
evaluate one algebraically and graphically. One should also understand how tell whether
or not a function is continuous.
II.) Derivatives: One should know the limit definition of the derivative (difference
quotient) and the various rules for finding derivatives (i.e. the power rule, the product
rule, the quotient rule, and the chain rule). One should also know how to take the
derivative of all trigonometric functions, exponential functions, and logarithmic
functions. One should also understand implicit differentiation.
III.) Applications of Derivatives: It is expected that one knows how to find the slope of
the tangent line and the normal line to the curve. One must be able to compare the graph
of a function to the graph of its first and second derivatives and vice versa. One must be
able to find local minima and maxima as well as points of inflection by using derivatives.
One should understand the Mean Value Theorem and Rolles Theorem. One should be
able to do optimization, motion, and related rates problems.
IV.) Integrals: One must know the various rules for antidifferentiation (i.e. the power
rule, u-substitution, the natural logarithm, inverse trigonometric functions). One should
know how to approximate the area under a curve by using rectangles and the Trapezoidal
Rule and get the exact area under a curve and between curves by using the 2nd
Fundamental Theorem of Calculus. One should also be familiar with the 1st Fundamental
Theorem of Calculus and accumulation functions.
VI.) Applications of Integrals: It is expected that one knows how to apply the integral to
problems in motion. One must know how to find the volume of a solid of revolution by
using the Disc Method, the Washer Method, the Shell Method, and various crosssections. One should know how to evaluate the average value of a function.
VII.) Differential Equations: One must be able to solve differential equations via
separation of variables. One should know how differential equations apply to
exponential change and Newtons Law of Cooling. One should also be able to
understand slope fields and linearization.
A Note on Technology:
Nowadays in advanced mathematics courses, the graphing calculator is used
extensively. While this machine is incredibly useful, it is important not to let it overcome
the power of the mind. Before using the calculator to solve a problem graphically or
analytically, the student should understand the theory behind the steps that the calculator
took to solve the problem. It is also important to remember that half of the AP Calculus
BC test prohibits the use of a calculator. Therefore, wherever possible, the student
should do calculus in the 19th-century fashion (even without a slide rule!) by hand. As
of the completion of the manuscript of this book, the TI-83 graphing calculators have
been discontinued. Thus, when discussing solutions via calculator, I will refer to its
contemporary counterpart, the TI-84. I would also like to mention that throughout the
book, notably in the final chapter (Vector Calculus), I have generated some of the graphs
using the computer program Mathematica. This program is one of the most powerful
computational tool in the world, and I encourage readers to learn more about its nearly
limitless capabilities.
Table of Contents
Prelude: At the Level of the Infinitesimal.5-7
Chapter 1: LHpitals Rule and Advanced Techniques
of Integration.8-24
Chapter 2: Differential Equations. 25-39
Chapter 3: Infinite Sequences and Series..40-55
Chapter 4: Power Series and Polynomial Approximations...56-74
Chapter 5: New Coordinate Systems: Parametric and Polar.75-102
Chapter 6: Vectors and Vector Calculus......103-131
Practice Test with Answers..132-188
Appendix A (Essential Pre-Calculus Information)..189-190
Appendix B (Brief Table of Integrals)..191
References.192
Index193-195
for the latter. Therefore, there must be a solution to Zenos paradoxes. This is where
calculus makes its entrance. While the mathematical philosophy of moving to an
infinitesimal level was manifest in antiquity, calculus did not emerge as a serious
discipline and practical tool until the early 18th century, when scientists were rigorously
attempting to find solutions to difficult problems in physics. Two prominent figures
during this period, Isaac Newton and Gottfried Leibniz, are usually credited with the
invention of differential and integral calculus. While there are some philosophical
naysayers who believe otherwise, calculus very concisely solves Zenos paradoxes.
Through calculus, it can be shown that even an infinite number of distances can have a
finite sum. Furthermore, it is not the case that an infinite span of time would be required
to travel an infinite number of infinitely small distances. Infinitesimal distance (dx) and
infinitesimal time (dt) have finite significance in differentiation and integration. For
dx
instance, the expression
= v , which means that the derivative of the position function
dt
is equal to the velocity function assumes a very helpful form when integrated:
dx = vdt , which means that an infinite sum of infinitely small distances yields the
infinite sum of products of instantaneous velocity and infinitely small time spans. This
relates infinitely small changes in position with infinitely small changes in time, so,
despite the notion of the infinitesimal, motion does exist!
It is often necessary in the sciences and business to move to the infinitesimal, or
differential, level. Why? Without doing so, one is restricted to differences between two
points and ignores what is in between! In business, for example, it is necessary to
optimize the volume of production to maximize profit. This is achieved by finding the
number of products at which the difference between revenue and cost of manufacture
(profit) is a maximum. Two isolated points for revenue and cost is not sufficient; an
infinite number of such points is needed! One must evaluate the derivative of the profit
function and find the point at which it undergoes a transition from positive to negative
values, a method covered in AP Calculus AB. In the pure and applied sciences, moving
to the differential level is indispensable in modeling phenomena. Similar to the business
example, calculus allows one to consider all of the points that model an event, not just
two isolated points. Furthermore, analysis of a differential aspect of a discrete body
easily allows for the extension to the whole body; that is, through the methods of
calculus, the mathematical equations that describe an infinitely small sliver of an object
also describe that object as a whole. This somewhat bottom-up approach is very often
more practical that analyzing the object as a whole. In physics, for instance, objects have
a property known as the center of mass, the point at which the mass of the object can be
considered concentrated. In effect, one can consider infinitely small chunks of an
object, each having a mass mi. When the average of these masses is found and weighted
(statistically) by each masss position with respect to a reference point, the center of mass
i=n
m r
Mathematically, center of mass = lim
m
i =1
i i
rdm
is found.
In this book, the theme of moving to the differential level will continue to be
emphasized as a useful method for understanding both the pure-mathematical aspect of
AP Calculus BC topics and their applications to the sciences.
x a
Firstly, the rule can only be used for limits in the indeterminate forms of and ; one
0
must not attempt to use LHpitals Rule to evaluate limits that are not in indeterminate
form. Thus, it is imperative that limits be first evaluated by conventional methods and
1 1
then, only if the need arises, be subject to the rule. Note also that the forms , , and
0
are not indeterminate; they have meaningful arithmetic significance. The first
0
fraction yields infinity, the second fraction yields zero, and the third fraction yields
infinity.
0
results,
0
the derivative of the numerator and denominator may be taken until a determinate form is
achieved.
The Indeterminate Form 0/0: If one attempts to take the limit of a quotient and
tan x
.
x 0
x
tan x tan(0) 0
= Use LHpitals Rule
Solution: lim
=
x 0
x
(0)
0
sec 2 x
sec 2 (0)
tan x
= lim
= lim
= 1.
x 0
x 0
x 0
1
1
x
lim
ln x
.
x x 3
ln x ln()
=
= Use LHpitals Rule.
x3
( ) 3
1
1
1
ln x
lim 3 = lim x 2 = lim 3 =
= 0.
x x
x 3 x
x 3 x
3() 3
Solution: lim
x
In the introductory discussion of LHpitals Rule, it was stated that the rule can
0
only be used to evaluate limits of the forms and . What of the other five
0
0 0
0 = 0 =
and 0 =
= .
1 0
Note that evaluating the product of two numbers and dividing one number by the
reciprocal of the other number are different operations that mean the same thing
3 2
arithmetically! For instance, 2 3 = = = 6 . Thus, while the above technique does
1 1
2 3
change the representation of the indeterminate form 0, it does not change the
arithmetic result.
Ex.) Evaluate lim x 2 e x .
x
Solution: lim x 2 e x = () 2 e ( ) = 0
x
x2
x 2 ( ) 2
Take the reciprocal: lim x e = lim
= lim x = ( ) = Use LHpitals
x
x 1
x e
e
x
e
2
2 x 2()
x
Rule: lim x = lim x = ( ) = Use LHpitals Rule again:
x e
x e
e
2x
2
2
2
lim x = lim x = ( ) = = 0 .
x e
x e
e
2
The Indeterminate Forms 1, 00, and 0: Since these indeterminate forms involve
exponents, a logarithm must somehow be applied for a reversal. If one takes the
natural logarithm of these forms, one can then deal with the problem in the context of the
techniques already discussed. Recall the property of logarithms that ln a b = b ln a . Once
the value of the limit is determined after taking the natural logarithm and applying the
previous techniques, it is important to remember to exponentiate (take the e of) of the
answer to know the original functions behavior.
1
x 1
(Note that the + subscript denotes a right-hand limit, which, along with lefthand limits, is discussed in AP Calculus AB. This is a necessary specification, since one
cannot approach a value of 1 from the left because the function does not exist to the left!)
lim x
1
x 1
x 1+
= (1)
1
(1) 1
1
1
ln(1) = 0 Take the reciprocal:
lim+ x x 1 = lim+
ln x =
x 1 x 1
x 1
(1) 1
1
1
(1) 1
x 1
1
ln x
ln(1)
1
x 1 ( x 1) (1 1) 0
=
=
= Use LHpitals Rule again:
lim
x 1+
ln(1) 0
1
ln x
ln x
1
x
10
1
1
1
1 1
= lim+
=
=
=
= . Note
x 1
0
0
1
x 1 ln x ln(1)
x ln x
2 2
1
x (1)
x 2
that is not the final answer. Since the problem was manipulated through the use of
the natural logarithm, this must be undone through exponentiation. Thus, the final
answer is e = 0.
The Indeterminate Form : Transforming this form into one that is workable
requires an algebraic manipulation of terms to yield two terms such that either one or
both are fractions. Once this is accomplished, one can multiply each term by a common
denominator to trivialize the subtraction. Since this explanation is undoubtedly not as
clear as the others, it is important to analyze the following example closely..
1
1
Ex.) Evaluate lim
x 0 sin x
x
1
x sin x (0) sin(0) 0
1
lim
=
= Use LHpitals Rule:
= lim
x 0 sin x
x x 0 x sin x
(0) sin(0)
0
x sin x
1 cos x
1 cos(0)
0
lim
= lim
= lim
= Use LHpitals
x 0 x sin x
x 0 x cos x + sin x
x 0 (0) cos(0) + sin(0)
0
Rule again:
1 cos x
sin x
sin(0)
0
lim
= lim
=
= = 0.
x 0 x cos x + sin x
x 0 x sin x + cos x + cos x
(0) sin(0) + 2 cos(0) 2
lim+
( x 1)
= lim+
x 1
ln x
11
d
dv
du
uv = u + v
d (uv) = udv + vdu
dx
dx
dx
The last equation can be rearranged into the form: udv = uv vdu . This equation
allows one to decompose an integral into two manageable parts. The most renown
example of the necessity for integration by parts lies in the integration of the natural
logarithm.
Ex.) Evaluate
ln xdx .
The form of this integral is deceptively simple. While it may seem that the
integration of this function is relatively straightforward, the techniques of AP Calculus
AB cannot be used here. Instead one must use the technique of integration by parts. One
dx
du
(because
, in
chooses u to represent ln x and dv to represent dx. In this way, du =
x
dx
this case, is the derivative of ln x) and v = x. Now the formula for integration by parts
dx
= (ln x)( x) dx = x ln x 1 + C .
can be applied: ln xdx = (ln x)( x) x
x
In this last problem, it was relatively simple to determine what to choose as u and
what to choose as v. This determination, however, is not always so clear. One should not
blindly choose these values, for making an incorrect choice could cost a great deal of
time and energy. While it may be intuitive to the reader which functions to choose, it is
easier to follow a heuristic (rule of thumb), conveyed by the acronym LIPET
(Logarithmic functions, Inverse trigonometric functions, Polynomials, Exponential
functions, and Trigonometric functions). Basically, if one encounters an integral that
must be solved via integration by parts, he or she should first look for a logarithmic
function to represent u. If a logarithmic function is not present, one should choose an
inverse trigonometric function to represent u, and so on, in accordance with the LIPET
algorithm. Also note that integration by parts may need to be performed several times
2 x 2 e 2 x dx .
12
e du = e
u
must be multiplied by 2 to obtain the form e u du , the integral must be multiplied by 1/2 in order to compensate for this algebraic manipulation. Again, this is AP Calculus
AB material).
1
1
1
1 2 x
2 x 2 e 2 x dx = (2 x 2 )( e 2 x ) e 2 x 4 xdx = (2 x 2 )( e 2 x ) +
e 4 xdx.
2
2
2
2
The integral is still not in a manageable form. Integration by parts must be
performed upon this last integral. In this case, u = 4 x , dv = e 2 x dx , du = 4dx , and
1
v = e 2 x .
2
2 x
4 xdx =
1
1
1
1
(4 x)( e 2 x ) e 2 x 4dx = (4 x)( e 2 x ) e 2 x 2dx = (4 x)( e 2 x ) e 2 x + C .
2
2
2
2
The last phrase is not the final answer. Recall that this is the solution to the
integral in the first round of integration by parts. That is, it is the solution to e 2 x 4 xdx .
1
The solution to this integral must be multiplied by 1/2 and added to (2 x 2 )( e 2 x ) . The
2
final answer is determined in the following way:
2 x 2 e 2 x dx =
1
1
1
1
1
= (2 x 2 )( e 2 x ) +
e 2 x 4 xdx = (2 x 2 )( e 2 x ) + (4 x) e 2 x e 2 x + C =
2
2
2
2
2
1
1
1
x 2 e 2 x + e 2 x e 2 x C = x 2 e 2 x + e 2 x C . This is the final answer.
2
2
2
Note that, since C is merely an unknown constant, 1/2C is merely another C.
Some integration by parts problems require addition of integral expressions, as the
following example shows:
Ex.) Evaluate
e cos xdx.
x
(sin x)(e )dx = (e )( cos x) ( cos x)(e )dx . Note that this last
x
mathematical phrase is the solution to the integral yielded by the first round of integration
13
Notice that the integral on the far right is the same as the integral that one wishes
to solve. This is where addition comes in:
e (sin x + cos x)
2 e cos xdx = (e )(sin x) + (e )(cos x) e cos xdx =
+C.
2
x
Integration by Parts Using the Tabular Method: While the previous problems were
undoubtedly mathematically tedious, imagine having to perform integration by parts
three, four, or more times for a single problem! There is, in fact, a method for integrating
by parts that is far more elegant (and easy!). This is known as the tabular method. The
following algorithm explains how to employ this useful technique:
1.) Make a table with three columns, labeled u, dv, and 1.
2.) Choose a value for u such that it will eventually yield zero when differentiated
enough (e.g. polynomials), put it in the first row of the u column and differentiate down
the column until zero is reached.
3.) Choose a value for dv and put it in the first row of the dv column and integrate
down the column. Constants of integration are not necessary here.
4.) Put a +1 in the first row of the 1 column and alternate between positive and
negative signs down the length of the column.
5.) Draw a diagonal line from each row in the u column such that it passes
through the appropriate row in the dv and 1 columns. Once the zero in the u column is
reached, cease drawing the diagonal.
6.) Multiply all three terms connected by a given diagonal and add all of these
multiplied terms to yield the final answer.
The algorithm sounds confusing, but it usually takes less than a minute to
compute the integral. The tabular method is best conveyed through an example:
Ex.) Evaluate
x 5 cos xdx. .
The polynomial ( x 5 ) will be chosen as u and cos xdx will be chosen as dv. With
these values, one may generate the table to solve the integral:
14
u
x5
5x 4
20x 3
60x 2
120 x
120
0
dv
cos x
sin x
cos x
sin x
cos x
sin x
cos x
1
+1
-1
+1
-1
+1
-1
+1
-1
Notice that the first diagonal connects the terms x 5 , sin x , and -1; the second diagonal
connects the terms 5 x 4 , cos x, and 1 ; the third diagonal connects the terms
20 x 3 , sin x, and +1; and so on. The terms of each diagonals are multiplied and these
multiplied terms are then added:
x cos xdx. =
5
x 5 sin x + 5 x 4 cos x 20 x 3 sin x 60 x 2 cos x + 120 x sin x + 120 cos x + C . The length of
this answer indicates that solving the problem without the use of the tabular method is
quite rigorous. Indeed, the calculation takes up a couple of pages if one uses the formula
for integration by parts alone! Note, however, that the tabular method can only be used
when the u term is not infinitely differentiable, i.e., when it is a polynomial.
Integration by Partial Fractions: Upon adding fractions with dissimilar
denominators, one looks for common denominators to yield one equivalent fraction. For
x
2x
(2 x)(4 x + 2) + (3x + 1)( x)
. This is elementary algebra.
instance,
+
=
3x + 1 4 x + 2
(3 x + 1)(4 x + 2)
Working backwards is slightly more difficult, and can often be much more difficult. How
can one decompose a fraction into a sum of fractions? To do this, one would use the
method of partial fraction decomposition. In this method, a fraction is decomposed
into a sum of fractions with denominators that are the factors of the denominator of the
original fraction and with numerators that contain unknown constants. For instance, the
x 2 + 3x + 2
, has a denominator with factors of x and x 2 + 2 x + 1 (i.e.
fraction, 3
2
x + 2x + x
3
x + 2 x + x = x( x 2 + 2 x + 1) ). Thus the decomposition is as
x 2 + 3x + 2
A
Bx + C
, where A, B, and C are unknown constants.
= + 2
3
2
x + 2x + x x x + 2x + 1
Notice how the degree of the numerator is one less than the degree of the denominator.
In the denominator of the first fraction, x is linear, so the numerator must be of zero
order, with a constant of A. In the denominator of the second fraction, x 2 + 2 x + 1 is
quadratic, so the numerator must be linear, with Bx+C. Note that it is possible for factors
( x + 2)
of the denominator to be repeated, as in the fraction
, which is decomposed as,
( x + 1) 5
follows:
15
( x + 2)
A
B
C
D
E
. Here, x+1 is known as the
=
+
+
+
+
5
2
3
4
( x + 1) ( x + 1)
( x + 1)
( x + 1)
( x + 1)
( x + 1) 5
repeated factor, and repeats five times in the decomposition process, with the value of
the exponent increasing from 5 to 1. Notice that the denominators are actually linear;
this is why all of the numerators are constants. The denominators are still linear despite
their being raised to a power. This is simply a repetition of linear factors. This
discussion of partial fractions could be proven by a corollary to the Fundamental
Theorem of Algebra. However, since this is not an advanced algebra textbook, it will not
be presented here.
Where does integration play a role in all of this? Similar to the philosophy of
integration by parts, certain expressions that are not integrable by conventional methods
can be broken down into more accessible components. In the method of integration by
partial fractions, a fraction that is not integrable through the use of the methods already
learned is decomposed into smaller fractions that can be integrated through the use of
conventional methods. Remember that these fractions contain unknown constants that
must be determined. This determination is first carried out by multiplying both sides of
the equation (i.e. the whole fraction side and the partial fractions side) by the
denominator of the whole fraction. This operation cancels the denominators on both
sides of the equation. When this is completed, one can solve for a system of equations.
x 2 + 3x + 2
dx .
Ex.) Evaluate
x3 + 2x 2 + x
x 2 + 3x + 2
A
Bx + C
= + 2
3
2
x + 2x + x x x + 2x + 1
Multiply both sides of the equation by x 3 + 2 x 2 + x :
x 2 + 3x + 2
Bx + C
A
= ( x 3 + 2 x 2 + x) + 2
( x 3 + 2 x 2 + x) 3
Remember that
2
x + 2x + x
x x + 2x + 1
x 3 + 2 x 2 + x = x( x 2 + 2 x + 1) .
x 2 + 3x + 2 = A( x 2 + 2 x + 1) + ( Bx + C ) x
x 2 + 3x + 2 = Ax 2 + 2 Ax + A + Bx 2 + Cx
x 2 + 3x + 2 = ( A + B) x 2 + (2 A + C ) x + A
This last equation is a standpoint from which one can solve a system of equations. Since
the coefficients on both sides of the equation must be equal,
( A + B) = 1
(2 A + C ) = 3
A=2
This simple system of equations is rather easy to solve. Since the constant A is now
known, this value can be substituted into the other equations that contain A to find values
for B and C. When this is done, it is determined that B = -1 and C = -1. With the
decomposition completed, the integral can be easily solved:
x 2 + 3x + 2
2
( x 1)
dx =
dx = 2 ln x 2 ln x 2 + 2 x + 1 + C .
dx +
3
2
2
x
x + 2x + x
x + 2x + 1
16
The system of equations in the previous problem was not at all daunting.
Oftentimes, however, the system of equations is so horrendous that one must use matrix
algebra to determine the unknown constants, which is not part of the AP Calculus BC
curriculum. In the field of mathematics known as linear algebra, matrices, arrays of
numbers that have a myriad of uses, can be used to solve a complex system of equations.
While a thorough discussion of matrix algebra will not be presented here, this point in the
book allows an excellent opportunity to use a graphing calculator to solve difficult
problems. In the case of a complex system of equations, the TI-84 graphing calculator
can represent a system of equations as a matrix and can then transform this matrix into
one of reduced row echelon form (RREF). A matrix in RREF has the following
properties: each non-zero row has a greater number of leading zeros than the previous
row, the first non-zero number in a row is 1, and the initial 1 of row is the only nonzero element in the column in which the 1 appears. Observe the matrix below to
clarify these requirements:
1 0 0 0 2
0 1 0 0 4
0 0 1 0 3
0 0 0 1 9
Notice how each row has one more has one more leading zero than the previous
one, though it is not a requirement that each row differ by just one zero. Note also that
whenever a value of 1 appears for the first time in a row, there are only zeros in the
column in which the 1 appears. Through the methods of linear algebra, namely a
method known as Gaussian elimination, a matrix that represents a system of equations
can be reduced to RREF so that the elements in the final column represent the values of
the unknowns. To understand how this is done, see the following example. Again, this
technique is not tested on the AP Calculus BC exam, but is very useful nonetheless.
3 x 4 + 4 x 3 + 16 x 2 + 20 x + 9
as the sum of three integrals.
Ex.) Represent
( x + 2)( x 2 + 3) 2
The first step is to decompose the fraction into a sum of partial fractions. Note
that the term x 2 + 3 is a repeated factor.
3x 4 + 4 x 3 + 16 x 2 + 20 x + 9
A
Bx + C
Dx + E
=
+ 2
+ 2
2
2
( x + 2) ( x + 3) ( x + 3) 2
( x + 2)( x + 3)
3x 4 + 4 x 3 + 16 x 2 + 20 x + 9
( x + 2)( x + 3)
=
( x + 2)( x 2 + 3) 2
A
Bx + C
Dx + E
=
+ 2
+ 2
( x + 2)( x 2 + 3) 2
( x + 2) ( x + 3) ( x + 3) 2
2
3x 4 + 4 x 3 + 16 x 2 + 20 x + 9 = A( x 2 + 3) 2 + ( Bx + C )( x + 2)( x 2 + 3) + ( Dx + E )( x + 2)
3x 4 + 4 x 3 + 16 x 2 + 20 x + 9 = Ax 4 + 6 Ax 2 + 9 A + Bx 4 + 3Bx 2 + 2 Bx 3 + 6 Bx + Cx 3 + 3Cx + 2Cx 2 + 6C
+ Dx 2 + 2 Dx + Ex + 2 E
17
3x 4 + 4 x 3 + 16 x 2 + 20 x + 9 = x 4 ( A + B) + x 3 (2 B + C ) + x 2 (6 A + 3B + 2C + D) + x(6 B + 3C + 2 D + E )
+ (9 A + 6C + 2 E )
A+ B = 3
2B + C = 4
6 A + 3B + 2C + D = 16
6 B + 3C + 2 D + E = 20
9 A + 6C + 2 E = 9
This is the system of equations that must be solved through matrix algebra. The
matrix will have as many rows as there are powers of x, including the power of zero for
the constant, and as many columns as there are unknowns plus one for the equality.
Since the system of equations in this example is associated with a fourth-order
polynomial, there will be five rows, and since there are five unknown constants, there
will be six columns. After the matrix is set up, each row is treated like one of the
equations in the system. For example, A + B = 3 , in matrix language, appears as
(1 1 0 0 0 3) , because the coefficient of A is one, the coefficient of B is one, there
are no Cs, Ds, or Es, and the whole equation is equal to 3. The final matrix is:
1 1 0 0 0 3
0 2 1 0 0 4
6 3 2 1 0 16
0 6 3 2 1 20
9 0 6 0 2 9
This matrix can be set up in the TI-84 calculator as follows: Go to 2ND + x-1,
which will open up the MATRIX menu. Go to EDIT and choose a matrix to edit (e.g.
[A], [B], etc.). The main screen of the calculator will ask for the number of rows and
columns. For this example, in put 56. Fill in the matrix.
This matrix must now be reduced to RREF. To do this via calculator, go to the
MATRIX menu again and go to the MATH option. Scroll down to B: rref( and select
this option. On the main screen of the calculator will appear the operation rref(. Input
the matrix (either [A], or [B], etc.) into this operator by choosing it from the MATRIX
menu. Press the ENTER key to calculate the RREF matrix. For this example, it should
1 0 0 0 0 1
0 1 0 0 0 2
be: 0 0 1 0 0 0
0 0 0 1 0 4
0 0 0 0 1 0
A =1
B=2
Therefore, C = 0
D=4
E=0
1
2x
4x
.
and the partial fractions are
+ 2
+ 2
( x + 2) ( x + 3) ( x + 3) 2
18
The problem merely asked for the setting up of the integral as a sum of integrals. Thus,
the final answer is:
3 x 4 + 4 x 3 + 16 x 2 + 20 x + 9
=
( x + 2)( x 2 + 3) 2
0.60976
1.19512 x + 1.60976
5.53659 x 3.07317
=
+
+
.
2
( x + 2)
( x + 3)
( x 2 + 3) 2
Improper Integrals
b
This theorem states that a definite integral is evaluated by finding the antiderivative of a
function, substituting the lower and upper limits of integration into this expression, and
subtracting the substitution of the upper limit from the substitution of the lower limit. In
the cases in which this theorem is employed, f(x) is bounded between a and b and the
function is continuous between and at these points (i.e. on a closed interval). A definite
integral with these properties is known as a proper integral. Whenever one or both of
these properties is not fulfilled, the expression is an improper integral. For instance, the
integral
+
2
Therefore, there are no restrictions on the interval. In another case, both limits may be
finite, but the function could exhibit a discontinuity within the interval in which
+1
1
dx is an improper integral
integration is to take place. For instance, the integral
1 x
because there exists an infinite discontinuity at x = 0. Interestingly, it may very well be
the case that an improper integral represents a finite area despite an infinite interval or an
infinite discontinuity. Such improper integrals are said to be convergent. Improper
integrals that do not represent a finite area are said to be divergent. In order to apply the
Second Fundamental Theorem of Calculus to improper integrals, it is necessary to define
the limit(s) of integration or the value of discontinuity that makes for impropriety as a
limit. This is the major theme of calculus yet again thinking on an infinitesimal level.
Indeed, it is not possible to think of an infinite geometric representation as having a finite
area. Nevertheless, by applying the definition of the limit to the problem, it is very much
possible!
1
from 1 x < finite? If so, calculate the
Ex.) Is the area under the curve y =
3x
value of this area.
+
1
dx
What the question is essentially asking is Does the integral
3x
1
converge?
To determine this, it is necessary to define the upper limit of integration as a limit so that
the Second Fundamental Theorem of Calculus can be used. Remember, infinity is not a
number, so it technically cannot be used in this theorem. Thus, the area (A) is defined as:
19
A = lim
A = lim
a +
a
1
1
dx , where a is some value that approaches infinity.
3x
1
1
1
dx = lim ln x 1a = lim [ln(a ) ln(1)] = 0 =
a + 1 3 x
3 a +
3 x+
Therefore, no finite area exists under this curve in the interval 1 x < .
3
5
5 5 x4
5 5 x4
a
1
A=
dx = lim
dx + lim
dx = 3 lim
1 + 3 lim+
b
5
5
5
4
4
a 0
b0
a 0
b 0
1 x
1 x
b x
4
5 4
5 5 (a ) 4 5 5 (1) 4
5
+ 3 lim 5 (1) 5 (b) =
= 3 lim
4
4
4
4
a 0
b 0 +
5 5
15 15
3 0 + 3 0 = +
= 0 . Note that while the integral equals zero, the area
4 4
4 4
3
15 15
below
is equal to 2 =
. To see why this is the case, observe the graph of y = 5
4
2
x
on the interval 1 x < +1 :
While one of the shaded regions would lie below the x-axis while the other would
lie above it, there is no geometric meaning to a negative area. Thus, the area under this
curve on the given interval is finite, with a value of 15/2.
Applications of Improper Integrals: While it may seem that the significance of improper
integrals is limited because they either involve limits of integration of infinity or points of
infinite discontinuity (or both), this is far from true. Before delving into the physical and
life science applications of improper integrals, it is important to consider them from
somewhat of a philosophical standpoint. This is where the strangeness of calculus comes
20
into play. While our intuition maintains that an infinite geometric figure cannot possibly
have a finite area or volume (after allit is infinite!), the methods of calculus beg to
differ. This is philosophically quite interesting. If one considers the non-Platonistic
premise that mathematics is a tool created by human beings, it is remarkable to note that
the human mind has birthed something that can explain something that it itself cannot
explain! Consider the case of Gabriels Horn, for instance. This geometric figure,
conceived by the Italian mathematician Evangelista Torricelli, has a finite volume, but an
infinite surface area. It can be modeled as the solid of revolution generated when a
1
hyperbola (e.g. y = ) is revolved about the x-axis:
x
Generated on Mathematica
To calculate the volume of this figure, on can use the methods of AP Calculus AB. In
this case, the disc method is the most appropriate:
2
a
1
1
dx
1
1
V =
= lim 1a = lim
+ = (0 + 1) = .
dx = lim
2
a ( a )
a 1 x
a x
(1)
1 x
Thus, the solid has a finite volume. But what of the surface area? While a topic of
neither the AP Calculus AB curriculum nor the AP Calculus BC curriculum, the surface
area of revolution (SA) for a curve f ( x) = y about the x-axis bounded on the closed
SA =
= 2 lim
dy
2y 1 + dx . For Gabriels Horn:
dx
2
1
1 + 2 dx = 2 lim
a
x
x
1
1+ 4
x
x
dx
1
1
1 1
ln x 2 1a = 2 lim ln a
+ 3 dx = 2 alim
2x
2(a) 2
x x
1
ln 1
2(1) 2
interesting scenario. If Gabriels horn were to hold paint, it would hold cubic units of
21
it. Obviously, this is not enough paint to cover its own surface, which is infinite! While
this rather philosophical problem has little practical significance (actually none), the
theory behind this problem plays an integral role (no pun intended) in the sciences in the
form of probability distributions. Note that the following section is not part of the AP
Calculus BC curriculum, but should be read to gain an appreciation of the material.
Probability Distributions in the Sciences: Very often in the physical and life sciences,
one must apply the laws of probability and statistics to solve problems. A probability
density function (often denoted as p(x)) is a mathematical representation that yields the
probability of the occurrence of a random variable x. For a probability distribution on
the interval < x < + , the improper integral is normalized, meaning that it is equal
to unity (or 1):
probability. Loosely speaking, if one considers all of the probabilities associated with a
probability density function, they must sum to 100 percent, or unity. Very often in
nature, the probability density function is a Gaussian function, or normal distribution
1 x 2
1
(a bell-shaped curve), represented as: p( x) =
exp
, where is the
2
2
mean and is the standard deviation. Observe the Gaussian function below:
While there are many variations on the Gaussian function (i.e. some may
skewed to one side, flattened, or stretched), it is imperative to note an important
characteristic of this function; it never equals zero. That is, no matter how far along the
x-axis, the probability of the random variables occurrence is always greater then zero, a
truth that has important implications in the sciences. Note the philosophical connection
to the improper integral problems already discussed; while the improper integral actually
represents a finite number (i.e. 1), the probability density function never touches the xaxis!
One of the most important applications of the probability density function to the
sciences is manifested at the molecular level. The physics applied to macroscopic
objects cannot be used to explain the behavior of the submicroscopic world. Rather, the
constituent components of matter must be described according to the laws of probability
and statistics. For example, consider a collection of gas molecules in a closed container.
If the gas is ideal, then it follows the postulates of the kinetic-molecular theory of gases.
This is a very important concept in the study of physical chemistry, for it describes the
submicroscopic behavior of matter. Where does the probability density function come
22
into play? Within the container described are many gas particles, perhaps on the order of
1023. The kinetic-molecular theory maintains that these particles move randomly in
straight lines and transfer energy and momentum during collisions but do not lose it.
All of the particles are moving at different speeds; one may be moving extremely
quickly, while another extremely slowly and these speeds change quite often.
However, at a given temperature there will be a most probable speed for the collection of
gas molecules. This is shown graphically as a Maxwell-Boltzmann distribution:
Note that the probability never decreases to zero as one moves farther along the x-axis.
Indeed, it is possible for a gas particle at 273 K (0C) to have a speed of, say, 1 10 50 m/s,
but it is not very probable! Notice the characteristics of the Maxwell-Boltzmann
distribution when the temperature increases:
The curve is flattening out such that the most probable speed is greater but is less
probable than the most probable speed at a lower temperature (i.e. the peak is lower).
Why is this case? Recall that for a probability density function the improper integral,
+
p( x)dx = 1 , holds true. Thus, in order to keep the area equal to unity when the most
probable speed is greater in magnitude, the probability of the most probable speed must
be lowered. What would the Maxwell-Boltzmann distribution look like at absolute zero
(0 K)? What about at a very, very high temperature?
Note that the Maxwell-Boltzmann distribution does not solely apply to gases, but
to all matter. This distribution is also important in understanding such processes as
diffusion, which is quite important to life on Earth.
The last application of the probability density function to be considered in this
chapter concerns the field of quantum mechanics. While this is, indeed, a broad field, it
isgenerally based upon the assumption that matter at the submicroscopic level, especially
electrons, exhibit the characteristics of waves more so than particles. As a result of these
wave characteristics, it is impossible to know with certainty both the position and
momentum of a particle simultaneously, a concept known as Heisenbergs uncertainty
23
While it is very unlikely that an electron will be located far away from the
nucleus, it is possible nonetheless. An electron associated with an atom in this sheet of
paper could be somewhere on the moon, but it is not very probable!
Concluding Remarks
This first chapter introduced several new techniques in differential and integral
calculus. LHpitals Rule, integration by parts, integration by partial fractions, and
improper integrals are all very useful tools when faced with more difficult problems that
cannot be solved through the use of the techniques of AP Calculus AB. While these
techniques are powerful, it is important to realize that they are not omnipotent. There are
many other techniques that can be used to evaluate troublesome limits and integrals, and
most are far beyond the scope of AP Calculus BC.
One very important application of improper integrals, the probability density
function, was also discussed in the context of its significance in modeling systems in
physical chemistry and quantum mechanics.
The next chapter will apply many of the techniques discussed in this chapter to
approach various scientific problems.
Key Terms:
indeterminate form
LHpitals Rule
integration by parts
LIPET
tabular method
partial fraction decomposition
repeated factor
integration by partial fractions
reduced row echelon form (RREF)
improper integral
convergent
divergent
Gabriels Horn
probability density function
random variable
normalized
Gaussian function
normal distribution
kinetic-molecular theory
Maxwell-Boltzmann distribution
Heisenberg uncertainty principle
wave function
24
Eulers Method
Eulers Method, pronounced oi-ler, is a numerical method used to approximate
solutions to differential equations. The method was developed by the Swiss
mathematician Leonhard Euler in 1768. Numerical methods for solving differential
25
equations are necessary when solutions cannot be found analytically, such as when one
does not explicitly know the algebraic structure of the differential equation, but knows
certain values for variables and slopes. Given the initial values of the variables and the
slope, a discretized (i.e. non-continuous) form of the limit definition of the derivative can
be used and rearranged to yield the formula for using Eulers method. Recall from AP
y ( x + h) y ( x )
. If h is not
Calculus AB the limit definition of the derivative: y ' ( x) = lim
h 0
h
y ( x + h) y ( x )
infinitely small (i.e. it has a finite value), the equation becomes: y ' ( x)
.
h
For the purposes of Eulers Method, let h be called x , let y ( x + x) be denoted as y n +1
and let y (x) be denoted as y n . The discretized form of the limit definition for the
y n +1 y n
. This is rearranged to yield the formula for
x
Eulers Method: y n +1 = y n + x y ' ( x) n . What is the significance of this formula? If one
knows initial values for a function and its derivative, while not necessarily knowing what
those functions are, and chooses a certain increment x , known as the step size, for the
numerical analysis, one can approximate the value of y(x) for a certain x. See the
example below.
Ex.) Approximate the value of y(1.5) by using increments of 0.1 if y(1) = 4 and
y(x) = x 2 y .
Notice that in this problem, the actual algebraic structure for the differential
equation is known. This sort of problem often appears on the AP exam, primarily as a
free-response problem. The rationale for a problem of this type will become clear at the
end of this example.
Since the starting x-value is 1 and the step size is 0.1, there must be six steps
involved to approximate a value for the function at x = 1.5. The first conditions given are
x0 = 1 and y 0 = 4 . Thus, the initial slope is (1) 2 (4) = 4 . The formula for Eulers
Method may now be used: y n +1 = y n + x y ' ( x) n y1 = (4) + (0.1)(4) = 4.4.
Five more steps to go! The next x-value is found by adding the step size to the previous
x-value: x1 = x0 + x = (1) + (0.1) = 1.1 . Now y1 can be found:
26
x
x +C
x
3
x3
2
x dx ln y = + C y = e
= e 3 (e C ) = Ce 3 .
C
Recall that C is merely an unknown constant, so e is just another C. This constant
can be determined based upon the initial values given:
dy
= x2 y
dx
y = Ce
x3
3
dy
=
y
(4) = Ce
(1) 3
3
C =
4
e
1/ 3
(1.5 ) 3
3
4
= 8.828287262
y = 1 / 3 e
e
This analytical solution differs from the numerical solution by 0.0837666646.
There are several key points to take away from this Eulers Method problem. Firstly, the
numerical method never yields exact solutions to differential equations. Secondly, as the
step size becomes infinitely small, the solution becomes exact. Thus, the smaller the step
size that one uses, the less error is involved in the calculation. However, notice that even
this problem with the relatively large step size of 0.1 was quite tedious to solve. One
must weigh the precision of the solution needed against the cost of actually solving the
equation numerically. Thirdly, Eulers Method is only one of many numerical methods
of solving first-order differential equations, and is generally the least powerful. In fact,
Eulers Method really only has historical significance, since no one uses this method any
more.
The Law of Exponential Change
Very often in the physical and life sciences, one encounters natural logarithms in
mathematical models of phenomena. This is the case because many aspects of nature
seem to change in direct proportion to an amount of something present. Whether this
amount refers to number of organisms in a population or the number of molecules in a
27
mixture of chemicals in a beaker, the change in these entities often takes on the form of
dN
the following differential equation:
= kN . This mathematical statement means that
dt
the change in some amount N over time is proportional to that amount. That is, if there is
more of something present, its rate of change is greater. To understand why so many
natural processes are modeled mathematically with natural logarithms (in fact, its natural
prevalence is one of the reasons for its name!), it is necessary to solve the differential
equation:
dN
dN
= kN
= kdt ln N = kt + C N = e ( kt +C ) = e kt e C = Ce kt .
N
dt
Let the initial value of N be N0: N 0 = Ce k ( 0) = C .
( )( )
28
carried out:
N
dN
1 1
1
1
kt + C =
=
= Kkt kC
+
dN = (ln N ln K N ) ln
N (K N ) K N K N
K
KN
(As always, kC will just be considered another C)
e Kkt
N
N
K
. This is logistic growth
= kC = Ce Kkt 1 = Ce Kkt N =
KN e
K
1 + Ce Kkt
equation. Note some important characteristics of its algebraic structure. The carrying
capacity is located in the numerator, the exponential is multiplied by a constant C, and e
is raised to the negative power of the product of the carrying capacity, another constant k,
and time. This structure puts a limit on the growth that the equation models. Compare
the exponential growth model to the logistic growth model:
Note that for the logistic curve g(t), lim g (t ) = K . The practical applications of
t
these curves will be discussed shortly. Note that logistic growth is frequently tested on
the AP exam either as a multiple-choice or free-response question.
The Learning Curve
In certain cases, the rate of change of a process is directly proportional to the
difference between the endpoint of a process and the degree to which the process has
already ensued. Mathematically, this is expressed as the differential equation:
dN
= k ( A N ), where A is the endpoint of the process. Suppose that the process in
dt
question is the typing of written statement on a word processor. Let the total number of
letters in the statement be 1,000 and let the efficiency be determined by the number of
correct letters typed relative to the 1,000 letters. If the process follows a learning curve,
or bounded growth, the efficiency will increase with time. What is the nature of this
increase? To answer this question, the differential equation must be solved:
dN
dN
dN
= k(A N)
= kdt
= kdt
dt
(A N)
(A N)
ln A N = kt + C ( A N ) = e kt C = (e kt )(e C ) = Ce kt
N = A Ce kt .
In the typewriting example, A would be 1,000 and C and k could be determined
from a set of initial conditions. Notice the graphical nature of the learning curve:
29
The rate of change is most rapid near the beginning and decreases to zero at N=A.
What does this mean? Basically, the greatest degree of learning takes place near the
beginning and decreases as the endpoint is reached. Thus, the more one learns (or the
more of a process that is carried out), the lesser the rate of increase of learning. Note also
that the efficiency, measured in this case as the difference between A and N, increases
with the time invested in the process.
The learning curve has many implications in economic and educational strategy.
If it is the case that as the amount of time invested in a project increases the efficiency
increases, then corporations and classrooms should concentrate on those strategies that
increase the experience of those involved. Problems concerning the learning curve
appear on both the AP Calculus AB and BC exams.
Ex.) An elementary school student must memorize the capitals of all fifty United
States in thirty minutes. The rate of memorization is directly proportional to the
difference between the number of capitals that must be memorized and the number of
capitals that have been memorized so far. Assuming that the child can flawlessly
memorize two state capitals during the first three minutes, and assuming that the child
knew none of the state capitals at the beginning of the study session, will he be able to
memorize all fifty state capitals in thirty minutes?
The problem did not require solving of the differential equation, so it would be
helpful in a problem like this to commit the resulting equation to memory. The same
applies to the exponential and logistic growth models as well, unless, of course, the
problem does ask for the derivation.
N = A Ce kt = 50 Ce kt
It is assumed that the child knew none of the state capitals before beginning
memorization, so (0) = 50 Ce ( 0 ) C = 50 .
ln(48 / 50)
0.0136073.
It is given that at t = 3, N = 2: (2) = 50 50e k ( 3) k =
3
Now, the value of N when t = 30 may be found:
N = 50 50e ( 0.0136073)(30) = 16.75834 17. Unfortunately, thirty minutes is not
enough for the youngster to memorize all of the state capitals.
Mathematical Models of Population Ecology
30
Ecology is a rather broad field of biology that studies the relationships between
organisms and their environment. One branch of this field, population ecology, focuses
on these relationships at the level of a group of the same species inhabiting the same area,
a group referred to as a population. Population ecology is perhaps the most quantitative
of the subfields of ecology, since it often studies the changes in the number of individuals
in a population over time, the strategies that different species use to proliferate, and how
biotic (living) and abiotic (non-living) factors affect a certain population. Populations are
generally modeled mathematically by the exponential growth model or the logistic
growth model.
The Exponential Growth Model in Population Ecology: The exponential growth model
for populations, specifically human populations, was devised by English economist
Thomas Malthus in his An Essay on the Principle of Population (1798). This treatise
exclaimed certain danger for the human race in that while the global food supply grows
linearly, the human population grows exponentially. Population ecologists model this
dN
= rmax N , which is essentially the same equation as the one
exponential growth as
dt
discussed earlier, but instead of using the constant k, rmax is used. This constant is known
as the intrinsic rate of increase, a measure of the capacity of a population to grow at its
maximum potential.
Ex.) A population of fruit flies (genus Drosophila) exhibits exponential growth.
There are initially 10 flies in the population. After three days, the population has
grown to 67 individuals. What is the intrinsic rate of increase of this population?
How many individuals will be present after one week has elapsed?
The problem did not require a derivation of the formula for the law of
exponential change, so it suffices just to commit the formula to memory. For an
example concerning population ecology, the formula is N = N 0 e rmaxt .
It is given that N 0 = 10 and that at t = 3, N = 67. From this information,
the intrinsic rate of increase can be determined:
ln(67 / 10)
(67) = (10)e rmax ( 3) rmax =
0.634036 .
3
Now the value of N when t = 7 can be found:
N = 10e ( 0.634036)( 7 ) = 846.268 846 flies!
At this point, it is probably a good idea to discuss the units in these problems.
The AP Calculus tests are usually not very strict in regards to including units in
calculations, as long as they are provided along with the final answer. Nevertheless,
especially when dealing with functions such as logarithms and exponents, it is important
to discuss the units of each element used in these mathematical models. In the formula,
N = N 0 e rmaxt , N and N0 have units of number of individuals, in the case of the previous
example, flies. The element t has units of time (seconds, minutes, hours, etc.), days in
the previous example. What about rmax? This element is not unitless. In the exponential
growth equation, it has units of reciprocal time, 1/days or days-1 in the previous example.
31
This is because exponentials; along with logarithms, trigonometric functions, and inverse
trigonometric functions; are known as transcendental functions, which are functions
that do not satisfy a polynomial expression. One of the consequences of the failure to
satisfy this expression is that the arguments of transcendental functions must be unitless.
Thus, in the example of the exponential growth model, the product of t and rmax must be
unitless. Since t has units of time, rmax must have units of reciprocal time.
The Logistic Growth Model in Population Ecology: While certain populations exhibit
exponential growth under certain conditions, no population can ever truly attain its
intrinsic rate of increase. A population could only truly grow in an exponential fashion if
it had unlimited access to resources such as food, water, and space. Obviously, these
resources are limited, as Malthus noted, meaning that there is a certain limit to the
amount of individuals in a population that an area can maintain. Indeed, if every
population underwent exponential growth, the total mass of existing organisms would far
exceed the mass of the Earth! Thus, the logistic growth model, discussed earlier in this
chapter, is a far more realistic way to gauge the trends in population growth. Population
dN
(K N )
ecologists often model logistic growth as:
= rmax N
, which is essentially the
dt
K
exponential growth model with another term added to tame the equation. This term,
(K N )
, represents the fraction of the population relative to the carrying capacity that is
K
still available for growth. Note that this is different from the logistic model already
discussed, in which the term was merely (K N), or the number of individuals that can
still undergo growth. While both differential equations yield a logistic function, the one
discussed first should be used for the purposes of AP Calculus BC.
Ex.) A population of grizzly bears in a preserve exhibits logistic growth defined
dN N
N
by the following differential equation:
= 1 . What is the carrying capacity of
dt
6 14
this population? How many grizzly bears are there when the rate of increase in the
population begin to decrease? If when t = 0 years, N = 2 grizzlies and when
t = 5 years, N = 10 grizzlies, determine the time at which this decrease occurs.
For the first part of the problem, recall that the carrying capacity is reached when
the rate of change of the population (the derivative) is equal to zero. Thus, the carrying
capacity can be found directly by setting the differential equation equal to zero are
N
N
solving for N: 1 = 0 N = 14 . This is the carrying capacity for the population.
6 14
Upon reading the second part of the problem, the reader may realize that this
question is essentially testing the ability to analyze the relationships between functions
and their derivatives, an AP Calculus AB skill. When the rate of increase (the derivative)
begins to decrease, the derivative will have a relative maximum. To determine the value
of N at which this occurs one could either graph the derivative on the TI-84 and
determine where its maximum lies, or use the second derivative test. To save time
(assuming that this is a question in which one is permitted to use a calculator!), one
should use the former method. Graphing the derivative, it is found that the derivative has
32
where the capital letters refer to particular chemical species and the lower-case letters
refer to their relative numbers in the reaction (called a stoichiometric coefficient). The
left side of the equation is referred to as the reactants and the right side of the equation is
referred to as the products. As the reaction takes place, the amount of reactants will
decrease and the amount of products will increase. How can one express the rate at which
this reaction takes place? In essence, one can take the derivative of the equation. There
is a problem, however; not all species involved may change at the same rate. For
instance, the coefficient b may be twice as great as coefficient a, meaning B will decrease
twice as fast as A. One can compensate for this by multiplying the rate of change of the
concentration of each species (concentrations are denoted with brackets [ ]) by the
inverse of its stoichiometric coefficient:
33
1 d [ A] 1 d [ B] 1 d [C ] 1 d [ D]
=
+
.
a dt
b dt
c dt
d dt
The negative signs on the left side of the equation signify a decrease in reactants
while the positive sign on the right signifies an increase in products. For instance,
consider the combustion of glucose ( C 6 H 12 O6 ) :
C 6 H 12 O6 + 6O2 6CO2 + 6 H 2 O.
The expression that conveys the reaction rate of this reaction would be:
d [C 6 H 12 O6 ] 1 d [O2 ] 1 d [CO2 ] 1 d [ H 2 O]
=
+
.
dt
6 dt
6 dt
6 dt
The actual nature of reaction rates can be explained by simple differential
equations that can be solved through separation of variables. Reaction rates are usually
gauged by the disappearance of reactants. Depending upon the chemical reaction in
question, these rates are dependant upon reactant concentrations in certain ways. For
instance, in a certain chemical reaction (AB), in which the species A is monitored, the
d [ A]
reaction rate might be expressed by the differential equation:
= k[ A], in which
dt
the rate of decrease of A is proportional to the amount of A present. This equation should
look familiar; its solution is the law of exponential change, but rather than expressing
growth, it expresses decay: [ A] = [ A]0 e kt . A reaction that proceeds in this manner is
said to be of first order because the rate of the reaction is directly proportional to the
concentration of reactants to the first power. In general, a reaction order refers to the
term n in the expression: rate = k [reactants] n , which is known as the rate law of the
reaction. If the rate of decrease of A were proportional to the square of the reactants, the
d [ A]
reaction would follow second-order kinetics (i.e.
= k[ A] 2 ). Reaction rates of
dt
higher order are certainly possible (some biochemical reaction orders are greater than
10), as are fractional orders, zero orders, and negative orders. All of these aspects of the
reaction are determined from laboratory experimentation.
Radiometric Dating
The atom is modeled as a collection of three subatomic particles: protons,
neutrons, and electrons. The electrons are distributed probabilistically (see chapter 1)
around an incredibly dense nucleus, which is composed of protons and neutrons. While
neutral (uncharged) atoms of the same element have the same number of protons and
electrons, they may very well differ in their number of neutrons. Atoms of the same
element that differ in their number of neutrons are referred to as isotopes. The stability
of an isotope is related to the ratio of its number of neutrons to its number of protons. An
unstable isotope is said to be radioactive, and will undergo a mode of radioactive decay
to achieve a more stable neutron-to-proton ratio. Scientists have learned to take
advantage of this submicroscopic phenomenon in the process of radiometric dating, the
determination of the approximate age of an object based upon the amount of a
radioisotope (i.e. radioactive isotope) that it contains, the known decay rate of the
radioisotope, and some reference object that also contains the radioisotope. One of the
34
most revolutionary forms of radiometric dating employs the isotope carbon-14. In this
technique, invented by Nobel-Prize winning chemist Willard Libby in 1949, one
compares the amount of carbon-14 in a sample (measured by the amount of radioactivity)
to the amount in living systems. This is done based upon the premise that the amount of
carbon-14 in carbon dioxide has been relatively constant for many thousands of years
and that the ratio of carbon dioxide with carbon-14 in its molecular structure to carbon
dioxide that does not contain radioactive carbon is the same in the atmosphere as it is in
living things. When a living thing dies, it ceases to assimilate carbon-14 into its
structure, and radioactive decay ensues. Thus, if one measures the amount of carbon-14
in an organically derived sample, such as something made from plant material, and
compares it to the amount of carbon-14 in a living thing, which represents how much
carbon-14 was originally present in the sample, he or she may determine the approximate
amount of time elapsed since the organism whose matter was used to make the object
died. The rate of decay of carbon-14 is not constant, but follows the law of exponential
change. This was determined by the observation that a radioisotopes half-life, the
amount of time elapsed before half of the material has decayed, is independent of any
external factors such as temperature or concentration. The formula for the half-life of an
object can be derived from the law of exponential change:
N = N 0 e kt When half of the material has decayed, half of the material
1
still remains, so N = N 0 . The time at which this occurs is the half-life, denoted
2
1
1
ln 2
as t1 / 2 N 0 = N 0 e kt ln = kt ln(1) ln(2) = kt t1 / 2 =
.
k
2
2
Notice that the half-life only depends upon the constant k, which is a
characteristic only of the specific isotope in question; other factors do not affect
half-life. This holds for other natural systems that follow the law of exponential
change as well. The half-life of carbon-14 is about 5700 years. Note that carbon
dating can only be applied to relatively young samples that are organic in origin.
After about 36,000 years, so much of the original radioactive sample has decayed
that it becomes quite difficult to detect, and the accuracy of the experiment is
substantially reduced. Problems concerning radiometric dating are common on
both the AP Calculus AB and BC exams.
Ex.) A medieval art collector plans to buy a flag that is claimed to have been used
in the Battle of Hastings in 1066. Before paying a hefty price of $50,000 for this
supposed relic, the collector demands that the seller have the flag carbon-dated. Analysis
shows that the radioactivity of carbon-14 detected in living plants from which such a
cloth may be made is 20.9 disintegrations per minute per gram and that the radioactivity
of the material in the flag is 18.6 disintegrations per minute per gram. Note that the halflife of carbon-14 is 5700 years. Is the flag authentic?
Let the variable A represent radioactivity. The radioactive decay of carbon-14
follows the law of exponential change: A = A0 e kt . While the initial radioactivity of the
sample (taken to be equivalent to the radioactivity of the living plant) and the final
35
radioactivity of the sample are known, the constant k is not. This constant may be found
ln 2
ln 2
from the formula for half-life: t1 / 2 =
k=
= 1.216048 10 4 .
k
5700
The approximate time elapsed may now be determined:
4
A = A0 e kt (18.6) = (20.9)e (1.21604810 ) t t = 958.7416 years . Since the Battle of
Hastings occurred 940 years ago, within a certain margin of error, the flag is probably
authentic.
Newtons Law of Cooling
Unlike the complex computational physics of today, the physics of Sir Isaac
Newton was known for its brevity and elegance. One example of this conciseness is
manifest in a simple law of heat transfer known as Newtons Law of Cooling, though the
model applies to heating as well. This model states that the rate of change of an objects
temperature is directly proportional to the difference between the objects temperature
and the temperature of the immediate environment, assuming that this environment is
large enough that it does not experience a substantial change in temperature.
dT
Mathematically,
= k (T Tenv ), where T is the temperature of the object at a certain
dt
time and Tenv is the temperature of the environment, which is assumed to be constant. If
an object is heating up, the value of k will be positive. If an object is cooling down, the
value of k will be negative. Note the important difference between this equation and the
equation for bounded growth; in Newtons Law of Cooling, the endpoint (the
temperature of the environment) is subtracted from the temperature of the object, while in
bounded growth, the variable in question is subtracted from the endpoint. This different
order makes a difference in terms of the final equation, as will be evident in the following
solution to the differential equation for Newtons Law of Cooling:
dT
dT
dT
= k (T Tenv )
= kdt
= kdt
(T Tenv )
(T Tenv )
dt
Note that this is a variation on the law of exponential growth; the only difference
here is that a constant (Tenv) is subtracted from the variable (T). What is the significance
of the constant C? This can be determined through an analysis of initial conditions. At
an initial time t = 0, C = T Tenv . This T is the initial temperature, or T0. Thus, the
equation for Newtons Law of Cooling becomes: (T Tenv ) = (T0 Tenv )e kt . Again, if the
object is heating up, k is found to be positive, and if it is cooling down, it is found to be
negative. Problems concerning Newtons Law of Cooling appear on both the AP
Calculus AB and BC exams.
Ex.) A blacksmith removes a piece of iron from a furnace with a temperature of
266C. The temperature of the blacksmiths workshop is 24C. When the iron has been
allowed to cool for one minute, its temperature is 235C. The blacksmith may work with
the iron when it cools to a temperature of 75C. Assuming that the rate at which the iron
cools is proportional to the difference between the temperature of the iron and the
36
temperature of the environment, how many minutes must the blacksmith wait until he can
handle the iron?
It is given that Tenv = 24 and that T0 = 266. To determine the value of k, one may
use the condition that when t = 1, T = 235 :
211
(T Tenv ) = (T0 Tenv )e kt (235 24) = (266 24)e k (1) k = ln
0.1370796
242
It can now be determined at what time the iron will cool to the desired
temperature:
ln (51 / 242)
11.35918.
0.1370796
Therefore, the blacksmith must wait approximately 11 minutes before handling the iron.
(T Tenv ) = (T0 Tenv )e kt (75 24) = (266 24)e ( 0.1370796 )t
Fd
Fg
37
From this free-body diagram, one can apply Newtons Second Law of Motion,
which states that the net vector sum of the forces acting on an object is equal to that
objects mass times its velocity: F = ma . In the scenario of an object falling through
the air with a relatively slow velocity: F = cv mg = ma , where m is the mass of the
object and g is the acceleration due to gravity, which is approximately 9.8 m/s2 . Recall
from AP Calculus AB that acceleration is the derivative of velocity. Thus, Newtons
dv
Second Law of Motion is essentially a differential equation: ma = m
= cv mg .
dt
While this differential equation is not tested on the AP exam, it still makes for a good
exercise in solving differential equations by the method of separation of variables while
also providing a meaningful application. This differential equation may be solved as
follows:
dv
mdv
mdv
m
= mg cv
= dt
= dt.
dt
mg cv
mg cv
These indefinite integrals can be made into definite integrals to simplify the
calculations. Assume that the initial velocity is zero (at time t = 0) and that any velocity
at time t is v. Then:
v
t
mdv
= dt = t t0 = t.
0 mg cv
0
The integral on the left can be solved through u-substitution. Let u = mg cv.
c
Then, du = -cdv, and the integrand must be multiplied by in order to yield the form
m
m
u / du = ln u + C , meaning the outside of the integral must be multiplied by c . Now:
m
m
v
t = ln mg cv 0 = [ln(mg cv) ln(mg )]
c
c
(mg cv)
mg cv
cv
mg
ct
= ln
e ct / m =
= 1
v=
(1 e ct / m ).
m
(mg )
mg
mg
c
This is a function of velocity with respect to time for the object as it falls through
the air. Notice the characteristics of this function when it is graphed:
38
isotopes
radioactive decay
radiometric dating
half-life
Newtons Law of Cooling
drag
free-body diagram
Newtons Second Law of Motion
terminal velocity
39
{ }
While one cannot evaluate limits of infinite sequences in the same manner as
functions, one can relate a certain infinite sequence to a function through this theorem,
and by determining the limit of the function, one can determine the limit of the sequence.
ln n
Ex.) Does the infinite sequence 3 have a limit as n approaches infinity? If
n
so, what is this limit?
The question is asking if the sequence converges or diverges. If the limit
exists, the sequence converges. If the limit does not exist, it diverges. In order to
40
determine this limit, it is necessary to employ the theorem that relates infinite sequences
ln x ln()
= . This limit is in indeterminate form. Therefore,
to functions: lim 3 =
3
x x
( )
ln x
1/ x
1 1 1
= = 0.
LHpitals Rule must be used: lim 3 = lim 2 = lim 2 =
2
x x
x 3 x x 3 x 3()
ln n
Thus, lim 3 = 0.
n n
This answer could also have been determined through the use of the TI-84. If one
presses the MODE key and chooses seq, he or she may graph sequences rather than
functions. Going to Y=, inserting 1 as nMin, typing ln(n)/n3, in the u(n) row,
and pressing the GRAPH key (one might need to choose an appropriate window for
viewing), one should notice two key features of the resulting graph. First, it seems to
approach the x-axis, which accords with the previous analytical solution of 0. Second,
the graph is not an unbroken curve, but a collection of dots. Since this is the graph of the
sequence, only certain, discrete values can be assumed, which are represented by the
dots:
Infinite Series
While sequences are merely lists of numbers, series are sums of the terms of a
sequence. In essence, infinite series and sequences are perhaps even more closely related
than this statement implies. Consider the sequence a1 , a 2 , a3 , a 4 ,L, a n . Suppose one
were to increasingly sum the terms of this sequence, i.e., start with a1 , then add a1 to a 2 ,
then add a1 , a 2 , and a3 , and so on. Each of these instances of increasingly encompassing
summation is called a partial sum, because, indeed, it is only a partial sum of the terms
of an infinite sequence. Thus, a new sequence is formed: S1 , S 2 , S 3 ,L, S n , in which
S1 = a1 , S 2 = a1 + a 2 , S 3 = a1 + a 2 + a3 , and S n = a1 + a 2 + a3 + L + a n . This new
sequence is an infinite series, an infinite sum of the terms of an infinite sequence. Series
are often not indicated as a sequence of partial sums, but in sigma notation:
= a1 + a 2 + a3 + L + a n .
n =1
41
At this point, a certain aspect of infinite series arises that defies intuition, but
which, nonetheless has a very real basis in calculus. While infinite series are the sum of
an infinite number of terms, this sum can be a finite number. Recall the discussion of
Zenos paradoxes from the prelude. The theory behind infinite series can actually solve
many of these paradoxes quite elegantly, as will become evident soon in this chapter. An
infinite series whose terms sum to a finite value is said to be convergent, while an infinite
series whose terms sum to infinity (i.e. a finite sum does not exist) is said to be divergent.
There are two ways in which an infinite series may converge. If a series
a is
n
n =1
absolutely convergent, then the sum of the absolute values of the terms is also
convergent, i.e.,
n =1
a converges, but a
n
n =1
does
n =1
n =1
become more meaningful as the chapter progresses, particularly because while the
convergence of a conditionally convergent series depends upon the order in which the
terms are summed, this makes no difference for an absolutely convergent series. Lastly,
there are three helpful algebraic properties of infinite series that are exactly analogous to
the algebraic properties of integrals:
1.)
(a + b ) = a + b
n
n =1
2.)
(a
n =1
n =1
3.)
bn ) =
n =1
a b
n
n =1
n =1
ca = c a .
n
n =1
n =1
a , lim a
n
n =1
diverges. However, if lim a n = 0, this does not prove convergence; it means that another
n
42
test must be used. It is easy to understand this from a logical standpoint; if the statement
If the limit does not equal zero, then the series diverges, is true, the reverse statement
If the limit equals zero, then the series converges, is not necessarily true. Only the
contrapositive of the original statement is definitely true: If the series converges, then
the limit equals zero. However, this statement cannot be used in the nth term divergence
test because the flow of logic must move from the truth about the limit to the conclusion;
the logic of the contrapositive statement moves from the conclusion to the truth about the
limit. If one knew the conclusion beforehand, one would already have the answer!
As a precaution, it is important to note that the limits that are evaluated in these
tests do not represent the actual sums of the series; their values merely convey whether
the series in question converges or diverges.
Test for Geometric Series: A geometric series is a series in which all terms have
a common factor and a ratio r that is raised to a specific power. Mathematically, a
ar
n =0
bound limit in the summation symbol is 0. This is often the case so that the common
factor a can be found easily, since it would be the first term in the sequence. There is
another way to make for this simplicity; a geometric series may sometimes be written as
ar
n 1
n =1
n
1
1 1
1 1 1 1 1
+L+
. Since
paradox is thus:
= + + +
2 2n
2 2 2 4 16 32
n =0
1
1
r = 0 < < 1 , the series does have a finite sum. Using the formula for the nth
2
2
43
a
(1 / 2)
=
= 1. Thus,
1 r 1 (1 / 2)
motion is not futile as the paradox would purport because the infinite number of distances
actually sum to a whole!
Test for P Series: A p-series is an infinite series in which the terms to be added
partial sum to find what this finite sum actually is: S n =
n1
n =1
1
1
1
1
+ p + p + L + p . A p-series
p
1
2
3
n
will converge if p > 1 , and will diverge if p 1. An interesting case that lies right on the
border of these conditions is the harmonic series, a p-series with p = 1, and, thus,
divergent. This name is derived from the physics of a vibrating string, which vibrates at
integer multiples of a fundamental frequency. These discrete multiples follow a p-series
in which p = 1. This is the mathematical basis for harmony in music. When one plucks
or bows a string instrument, or plays a wind instrument for that matter, it is not the case
that only one frequency is generated; most music would sound awful if this were how
acoustics worked! Instead, one actually perceives an infinite sum of frequencies that
follow a harmonic series (mostly some instruments do deviate slightly from the
harmonic series).
Telescoping Series: Besides the geometric series, there is one other series whose
exact sum can be very easily found the telescoping series, in which all but one or a few
terms will cancel each other through subtraction. Telescoping series are given their
informal name based upon the fact that they collapse like a pocket telescope, meaning
that most of the terms of this series disappear, leaving only a remnant of the fully
expanded form of the series. Telescoping series contain an infinite number of terms and
their opposites, so one can often spot such series by the presence of subtracted terms. See
the example concerning telescoping series at the end of this discussion of tests for
convergence and divergence.
The Integral Test: Recall from the beginning of this chapter that it is not possible
to perform the same calculus on discrete expressions like sequences and series as on
functions. However, there are mathematical theorems that allow one to analyze
discretized expressions as if they were functions. One such theorem is the integral test,
n =1
continually decreases, with positive terms, and that f(x) is also a monotonic decreasing,
positive function, which is also continuous, such that f (n) = a n for every positive integer
n =1
f ( x)dx
1
employ the ratio test. The theorem of this test states that a series
a converges when
n
n =1
44
a n +1
a
< 1 and diverges when lim n +1 > 1. However, if the limit is equal to 1, the test
n a
n a
n
n
is inconclusive. This test is particularly helpful in areas such as computer science and
combinatorics, a branch of discrete mathematics that studies (as its name suggests)
combinations of entities. This is because the factorial is an operation that expresses the
different ways in which entities may be arranged.
The Root Test: While the root test is rarely considered on the AP Calculus BC
exam, it is quite powerful. It is very closely related to the ratio test and is often
considered more powerful than the ratio test, because where the ratio test fails the root
lim
test often succeeds. The theorem of the root test states that a series
a converges if
n
n =1
lim n a n < 1 and diverges when lim n a n > 1. Similar to the ratio test, if the limit is
n
equal to 1, the test is inconclusive. Thus, since this makes the series appear smaller
than it would by using the ratio test, a series that appears too large to be convergent in
the ratio test might actually be proved convergent in the root test.
The Comparison Test: If one considers two infinite series, one which is known to
be convergent or divergent and one which is not, he or she can determine if the
unknown series converges or diverges by comparing it to the known series. If there
n =1
n =1
n =1
n =1
holds true:
n =1
n =1
converge if the larger series converges, but not necessarily vice versa), and
n =1
n =1
(2
n =1
+ 1 and
(2
n =1
the second series will have larger terms. This might not be so clear on other occasions.
The Limit Comparison Test: Oftentimes, one is not completely sure that every
term in one series is greater than the corresponding term in another series. However, the
45
series can still be compared through the limit comparison test. The theorem of this test
n =1
n =1
an
= L , where L 0 and L , then both series behave in the same way, i.e. they
n b
n
either both converge or both diverge.
The Alternating Series Test: Besides the telescoping series, which could have
negative terms, all of the series discussed so far have had positive terms. In an
alternating series, the terms oscillate between positive and negative values, making
these series suitable for modeling oscillatory phenomena. Alternating series are
characterized by the presence of 1 raised to some power. For instance,
lim
(1) a , (1)
n
n =1
n 1
a n , or
n =1
(1)
n +1
n =1
converge, three conditions must be fulfilled: all a n s are positive, a n > a n +1 for all n, and
lim a n = 0 . Do not become confused by the first part of the theorem; while the a n s
n
cannot be negative, the terms most certainly can, and will be! The negative 1 to an
exponent is not part of the a n .
It is now appropriate to bring up the topic of absolute and conditional
is convergent
n =1
n =1
n =1
if
n =1
n =1
56n
converge or diverge?
n =1
46
6
5
th
term
n =1
divergence test; if this test proves that the series diverges, one need go no further.
Remember, however, that this test does not prove convergence. Applying this test on this
1
1
series, lim 3 =
= 0. Thus, the nth term divergence test is inconclusive; another test
2
n n
( )
is needed. This series is easily recognizable as a p-series with p = 3. Since 3>1, the series
converges.
Note that, with enough mathematical experience, one could recognize this
series as a p-series and skip the nth term divergence test altogether. However, the next
problem will show that it can be quite useful for more algebraically complicated series.
n =1
3n 2 + 6n 1
converge or diverge?
n2
While more sophisticated tests could be used, one would benefit from first
th
applying the n term divergence to a complicated a n such as this:
3n 2 + 6n 1
lim
= 3. Thus, since this number does not equal zero, the
n
n2
series diverges. This limit was determined through the use of a theorem regarding limits
approaching infinity from AP Calculus AB. Recall the theorem that states that a limit
approaching infinity of a ratio of algebraic statements that have the same order is equal to
the ratio of the leading coefficients (i.e. the numbers multiplied by the variable raised to
the highest power) in the numerator and denominator, in this case, 3 and 1.
9 27 81
+
+ + L. Does this
2 4
8
geometric series converge or diverge? If it converges, what is its sum?
Ex.) The following is a geometric series: 2 + 3 +
n
3
3
geometric series can be represented as
2 . Since > 1, the series diverges.
2
2
n =0
4 +14n
converge or diverge?
n =1
Trying the nth term divergence test results in a limit of 0. Thus, another
test may be used. This is a power series with an extra term added in the denominator.
While other tests could be used, the comparison test will yield an answer more quickly.
47
Let
n =1
1
be the smaller series and let
4 + 4n 2
n =1
1
be the larger series. One can
4n 2
easily determine if the larger series diverges because it is a simple p-series. Since p = 2,
which is greater than 1, the larger series converges. Thus, according to the theorem for
4 +14n
n =1
1n n +1 4 .
n =1
There are no special formulae to determine the answer; the first several
terms must be written out and one must decide which terms cancel and which remain:
1 1
, , and, if the series were to be written out in its entirety, an infinite
5 6
1 1
1
number of fractions after that! The only terms that do not cancel are 1, , , and .
2 3
4
1 1 1 25
Thus, the sum of this telescoping series is: 1 + + + = .
2 3 4 12
n =1
n
converge or diverge? If it converges, does it
2e n
x
improper integral
dx and rearranges it algebraically for ease of analysis:
x
1 2e
a
1 x
1
xe dx = lim
xe x dx. This integral must be solved through the method of
2 1
2 a 1
integration by parts. One notices that one of the terms is a polynomial (x), so the tabular
method can be used:
48
u
x
1
lim
2 a
xe x dx =
dv
1
+1
ex
-1
ex
+1
-1
1
lim xe x e x
a
2
a
1
{[
][
]}
1
12 1
lim (a )e ( a ) e ( a ) (1)e (1) e (1) = = .
2 a
2e e
It must be determined if
n =1
n
converges or diverges. For this, one could use the
2e n
n =1
n =1
n
, which is divergent
1
n
for all positive n. This comparison series is divergent by the
2e n
nth term divergence test. Thus, the comparison test is inconclusive and one should move
on to the limit comparison test:
n
1
1
1
2e n
= lim n =
=
= 0. Thus, the series do not behave in
lim
(
)
n
n 2e
n
2e
1
n =1
n
, must converge. Therefore,
2e n
n =1
n
is
2e n
absolutely convergent.
10n!
4n
converge or diverge?
n =1
Since this series contains a factorial, the ratio test is probably the best
choice here:
(n + 1)!
4( n+1)
4n
10
= lim (n + 1)! 10 Note the property of factorials that
lim
4( n+1)
n!
n
n 10
n!
4n
10
(n + 1)! 10 4 n
10 4 n (n + 1)n!
(n + 1)
=
= lim
(n + 1)!= (n + 1)n!. Thus, lim 4( n +1)
. The
.
lim
n
+
4
(
1
)
n 10
n (10 4 )
n! n 10
n!
algebra was a bit cumbersome to yield this last fraction. While the n!s cancelled well,
( )
49
ab
(n + 1)
= a b c . That aside, lim
= . Since infinity is greater than 1, the series
c
n (10 4 )
a
diverges.
and
n
2
n =1
+1
converge or diverge?
The nth term divergence test is of no use here because the limit is equal to
zero. This is neither a geometric series nor a p-series. While other tests could be used on
this series, the least cumbersome is probably the limit comparison test. This series can be
n =1
n2 +1
n2 n
n5/ 2
. Since this is a
lim
= lim 2
= lim 5 / 2
1/ 2
n n
n n
n n
n
n
n
+
+
n2
limit involving infinity in which the orders (i.e. highest powers) of the numerator and
denominator are the same, one takes the ratio of the leading coefficients, which is 1 in
this case. Thus, the limit is exists and is non-zero. Therefore, both series behave in the
n
n =1
n
2
+1
to
converges.
n 2 2n 1
2
converge or diverge?
Ex.) Does the series
+
5
n
16
n
12
n =1
This series would be easy to analyze if the n power were not present. A
good rule-of-thumb is, if the series has at least one term raised to the n with no factorials
present, the root test is probably the best bet. This is the test that will be used in this
problem:
n
n 2 2n 1
n 2 2n 1
= lim 2
lim n 2
. Notice how this gets rid of
n
n 5n + 16n 12
5n + 16n 12
the troublesome n exponent. Since this is now a limit involving infinity with the same
order in the numerator and denominator, one can take the ratio of the leading coefficients:
n 2 2n 1
1
1
lim 2
= . Since < 1, the series converges.
n 5n + 16n 12
5
5
50
n =1
n
converge or diverge?
(5) n 1
n 1
n
1
(1) n 1 n 1 .
n =
5
5
n =1
n =1
n =1
Now one can decide whether this series is convergent or divergent by
using the three criteria described for an alternating series:
Every a n is positive.
n
=
(5) n 1
?
?
5 n ? (n + 1)
(n + 1)
n
n 1
(
5
)(
n
)
>
(
5
)(
n
+
1
)
>
5
n
>
(n + 1)
n
5 n 1
5 n 1
5n
?
? 1
1
4n > 1 n > . Yes, if n 1 , every n will be greater than .
4
4
( )
n
5
1
1
1
n
lim n 1 = lim n 1
= ( ) 1
= = 0.
n 5
n 5
ln 5 5
ln 5
a n > a n +1
>
51
n =0
a
of this series can be determined analytically: S n =
=
1 r
1
2 . The exact sum
2
1
2 such that the first five terms
series converges. Consider truncating the series
2
n =0
are summed, which is referred to as the fifth partial sum. How close is this value to the
1 1 1
actual value of 4? S 5 = 2 + 1 + + + = 3.875. While this seems rather close to the
2 4 8
actual value, it is helpful to define a parameter that will unequivocally indicate how far
off the estimate is. This parameter is known as the truncation error, which is, of course,
the error associated with splitting up the infinite series and summing a representative
number of terms. One can define the truncation in this context (Note that the term
truncation error will take on a slightly different meaning in the next chapter) to be the
percent error between the exact sum and the estimated sum. For the fifth partial sum:
4 3.875
error =
100 = 3.125%. Relatively speaking, this is a good estimate. While
4
this example conveyed the theory behind truncation, it was rather meaningless in a
practical sense since the exact sum of the geometric series was known. It is far more
interesting to consider examples in which only approximation works, but in which one
can also determine the truncation error.
The case of practical truncations to be discussed here is often tested on the AP
Calculus BC exam, and applies to alternating series. While the proof will not be shown
here, a theorem of alternating series states that the absolute value of the maximum
n =1
summing the first i terms, is less than or equal to the value of the (i + 1)st term: Ri ai +1 ,
where Ri is the truncation error associated with taking the sum of the first i terms.
Ex.) Calculate the error bound when the sum of the alternating series
th
n =1
(1) n n
10 n
52
While the problem implies that this alternating series converges, one must
always first prove that it does indeed converge using the alternating series test:
Every a n is positive.
n ? (n + 1)
>
Since an exponential function
10 n
10 n +1
increases much more rapidly than a linear function, the denominators of these
fractions will become larger more quickly than the numerators. Since the
n (n + 1)
(n + 1)
magnitude of the exponential for
is greater than that for n ,
will
n +1
10
10
10 n +1
be smaller than the other fraction because of its rapidly increasing exponential of
n
(n + 1)
higher magnitude in the denominator. Thus, n >
.
10
10 n +1
Prove that a n +1 < a n :
n 10 n
= 0 true.
(1)
n =1
n 1
n
5 n 1
to an accuracy of 5
decimal places?
This is the series from the example in the last section, which was proved to
converge. Thus, the proof need not be carried out here.
The problem asks for an accuracy of 5 decimal places, which means that
Ri 0.00001. From this information, one must determine how many terms (i) must be
summed to approximate the value within 5 decimal places:
j
Ri ai +1 a i +1 = 0.00001 Let i + 1 = j j 0.00001 . This is quite
5
difficult to solve analytically, because the variable is tied up both in a linear function and
in an exponential function. It is a good idea to solve this graphically using the TI-84.
x
One graphs y = x and y = 0.00001 and determines where they intersect. These two
5
graphs intersect at x = 8.4817448. Thus, (i + 1) must be less than or equal to 8.4817448.
The value i must be the next integer less than 8.4817448, which is 7. Therefore, by
summing the first 7 terms of the series, one achieves an accuracy within 5 decimal places:
53
(1)
n =1
n 1
n 1
(1) 0
7
6
5
4
3
2
1
+ (1)1 1 + (1) 2 2 + (1) 3 3 + (1) 4 4 + (1) 5 5 + (1) 6
0
5
5
5
5
5
5
5
A Case Study of Infinite Series Approximation: Will We Ever Know The Value of ? :
It is quite possibly the most enigmatic numerical expression of all human history.
For millennia, mathematicians and scientists have arduously striven to determine the
value of to more and more decimal places. Most of the numerical techniques
employed in these toils involve infinite series. While the Egyptians and Babylonians
were quite successful in their computations, extending the value of to several decimal
places by 1000 CE, it was arguably not until the late 14th century when the venerable
Indian mathematician Madhava of Sangamagrama, controversially considered by some
to be the true founder of calculus, developed a highly effective infinite series with which
he determined to 11 decimal places:
1
(1) n 1 12
1 1
= 12 1 +
9 45 189
n =1
another Indian mathematician, Srinivasa Ramanujan, developed another infinite series
that has spearheaded the modern computational advances in the calculation of .
Ramanujans infinite series, though enormously helpful to other mathematicians, is not
very elegant:
2 (1103 + 26390n)
29801
[(n!) (396) ] .
4
4n
n =0
history are manifested in the continual computational search for more and more digits to
.
Key Terms:
sequences
infinite sequences
series
partial sum
infinite series
absolutely convergent
conditionally convergent
nth term divergence test
geometric series
p-series
harmonic series
telescoping series
integral test
ratio test
combinatorics
root test
comparison test
limit comparison test
alternating series
truncation
truncation error
Madhava of Sangamagrama
Srinivasa Ramanujan
55
a ( x c) , where a is an
n
n =0
expression with ns, the same as the expressions associated with the series discussed in the
previous chapter, x is a variable, and c is a constant at which the series is centered. Note
that power series have lower limits of 0. Besides the fact that power series contain
variable terms (i.e. xs), there is yet another intricacy that must be introduced: It does not
suffice to specify that a power series converges; one must also specify where it
converges. A power series may converges in three ways: 1.) It could converge only at the
value c at which it is centered, 2.) It could converge at all real numbers within a certain
radius r (called the radius of convergence) from the center c (and, further, this could
include the value r of the radius, or it could not), or 3.) It could converge at all real
numbers. The complexity of the second case must be discussed further. Suppose the
graph below represents a power series centered at the origin (i.e. c = 0):
The radius to the right of c has a value of (c + r ) and the radius to the left of c has
a value of (c r ) . Note that c is positioned at the origin here for simplification; a power
56
series could have a non-zero value of c as well. If the power series represented in the
graph converges within the radius of convergence, it must be specified if this does or
does not include the values of (c + r ) and (c r ) . That is, it must be specified whether
the interval of convergence, all of the values at which the series converges, is a closed
interval or an open interval, an intricacy discussed in Calculus AB or, for those students
who have taken the class, pre-calculus. Note that while the term radius of convergence
is usually only applied to situations like that implied in the graph, it could have a value of
zero if it is only convergent at c or it could have a value of infinity if it is convergent at
all real numbers.
How can one determine the radius of convergence of a power series? In general,
one uses the ratio test introduced in the previous chapter because it gives concrete
a
numerical conditions for convergence; the series will converge when lim n +1 < 1 .
n a
n
Ex.) Determine the radius and interval of convergence of the power series
n =0
n 2 ( x 1) n
.
4n
a ( x c) . In this is example,
n
n =0
n
and c is equal to 1. One can use the ratio test to determine at
4n
what values of x this series converges:
the a n term is
(n + 1) 2 ( x 1) ( n +1)
4 ( n +1)
lim
n
n 2 ( x 1) n
4n
2
( n +1)
4n
= lim (n + 1) ( x 1)
n
4 ( n +1)
n 2 ( x 1) n
(n + 1) 2 ( x 1) (n 2 + 2n + 1)( x 1)
= lim
=
4n 2
4n 2
n
1
= ( x 1). Note that the (x 1) remains in the same form because it does not
4
contain any ns, so it is not part of the limit. The ratio test maintains that a series
a
1
converges when lim n +1 < 1 , so the series will converge when ( x 1) < 1 . From this
n a
4
n
expression, one can determine the radius and interval of convergence:
1
1 < ( x 1) < 1 4 < ( x 1) < 4 3 < x < 5 . While the radius of
4
convergence is definitely 4, one must test both endpoints of -3 and 5 to determine if the
interval of convergence includes these points, one must substitute them for x in the
original series:
n =0
n 2 (3 1) n
=
4n
n =0
n 2 (4) n
=
4n
n =0
57
n =0
n 2 (5 1)
=
4n
n =0
n 2 (4) n
=
4n
n =0
Thus, the interval does not include the endpoints, so the interval of convergence is
solely 3 < x < 5. Note that, in this case, the nth term divergence test was sufficient to
prove divergence when x was -3 and 5. However, for other problems the other tests
discussed in the previous chapter may be necessary.
n =0
n p xn
, determine the value, or values, of p such
3n
p ( n +1)
3( n +1)
x(n + 1) p 1
3n
= lim (n + 1) x
=
= x
lim
lim
n
n
3
n p x n n 3n p
3( n +1)
n p xn
n
3
In this power series, a n =
1
x < 1 3 < x < 3.
3
One must now determine for what value, of values, of p the following
series will converge:
n =0
n p (3) n
=
3n
n =0
and
n =0
n p (3) n
=
3n
n =0
58
represented as a geometric power series. Recall that the nth partial sum of a geometric
a
1
series is equal to
. Thus, for the function f ( x) =
, a = 1 and r = x. Thus, the
1 r
1 x
function may be expressed as follows:
1
f ( x) =
=
1 x
= 1 + x + x2 + x3 + L xn .
n =0
n =0
power series converges on a certain interval of convergence, the function that it defines
has an infinite number of derivatives inside this interval. If one differentiates each term
of the series, the resulting series will also converge within the same interval of
convergence, but not necessarily including the endpoints. The action of taking the
derivative of each term of such a power series is known as term-by-term
differentiation. This is easier to understand mathematically:
If f ( x) =
a n ( x c) , then f ( x) =
n
n =0
f ( x) =
na ( x c)
n
n 1
n =0
n(n 1)a ( x c)
n2
n =0
a ( x c)
n =0
series,
n =0
a n ( x c)
n +1
n +1
series, but not necessarily at the endpoints. The action of integrating each term of such a
power series is known as term-by-term integration. It is important to note a key
difference between the term-by-term integration theorem and the term-by-term
differentiation theorem; the term-by-term integration theorem does not state that one can
expect an endless number of integrals of a power series to converge within the same
interval convergence as that original power series only the first integral.
How are these theorems useful? Suppose one must know the power series
associated with ln(1 x) but only knows the geometric power series that defines the
x = 1+ x + x + x +Lx .
1
One could integrate this power series term-by-term because
1 xdx = ln(1 x) + C.
1
, which, as derived previously, is
function f ( x) =
1 x
n =0
However, one wishes to cause the constant of integration to disappear in deriving the
59
series for ln(1 x) . Thus, one can perform a definite integral and define f(x) in terms of a
new variable t:
ln(1 x) =
x
0
3
1
dt =
1 t
n =0
x n+1
t2 t3 t4
= t + + L
n +1
2 3 4
x2 x
x4
= x +
+
L
2
3
4
This last expression is the power series that represents ln(1 x). Thus, it is very
helpful to recognize a known power series as a derivative or integral of an unknown
power series. But how are these known power series determined in the first place?
This was a relatively simple task in regards to the geometric series because the function
a
was in the form of S n =
. How might one go about finding an explicit function for
1 r
(1) n x 2 n +1
? In order to simplify this problem, one wishes to
the power series
n
+
2
1
n =0
somehow create a geometric power series so that a function may be easily generated in
a
accordance with the formula S n =
. This is first carried out through term-by-term
1 r
differentiation:
f ( x) =
n =0
(1) n 1 (2n 1) x 2 n 2
=
2n 1
(1)
n 1
x 2 n 2 = 1 x 2 + x 4 x 6 + L.
n =0
(1) n x 2 n +1
is f ( x) = tan 1 x.
integration is 0. Thus, the explicit function for
2n + 1
n =0
2
60
Power Series
((21)n +x1)!
n
2 n +1
x
x
x
+
+L
3
!
5
!
7
!
n =0
Interval of Convergence: < x <
= x
x2 x4 x6
(1) n x 2 n
= 1
+
+L
(2n)!
2! 4! 6!
n =0
Interval of Convergence: < x <
n =0
x3 x5 x7
(1) n x 2 n +1
= x
+
+L
2n + 1
3
5
7
xn
x2 x3 x4
= 1+ x +
+
+
+L
n
!
2
!
3
!
4
!
n =0
Interval of convergence: < x <
n =0
(1) n ( x 1) n +1
( x 1) 2 ( x 1) 3 ( x 1) 4
= ( x 1)
+
+L
n +1
2
3
4
n =0
(1) n x n = 1 x + x 2 x 3 + x 4 + L
ex
ln x
1
1+ x
n =0
tan 1 x
cos x
Interval of Convergence: 1 x 1
Function
sin x
(1) n ( x 1) n = 1 ( x 1) + ( x 1) 2 ( x 1) 3 + ( x 1) 4 + L
1
x
An often tested concept on the AP Calculus BC exam is that a certain power series can be
manipulated algebraically to yield another power series. For instance, consider the last
1
power series in the table, for f ( x) = . One can easily find the power series for
x
1
f ( x) = 2 by simply substituting x 2 for every x. Thus,
x
1
=
(1) n ( x 2 1) n = 1 ( x 2 1) + ( x 2 1) 2 ( x 2 1) 3 + ( x 2 1) 4 + L. One can also
2
x
n =0
f ( x) =
n =0
x n f ( n ) (0)
x2
x3
= f (0) + xf (0) +
f (0) +
f (0) + L , in which the
n!
2!
3!
(0) indicates that the series is centered about the origin and the (n) denotes the nth
derivative of f (x). Such a power series is known as a Maclaurin series. While such
series had earlier been described by Madhava of Sangamagrama (see previous chapter)
and 17th century mathematician James Gregory, they are named for a later contributor,
the 18th-century Scottish mathematician Colin Maclaurin, who mentioned their use in his
Treatise on Fluxions (1742).
The above power series is influential because, in essence, it simplifies
complicated functions to polynomials of certain order; the higher the order of the
polynomial, the better the approximation to the complicated function in question.
Consider f ( x) = sin x as a simple instance of employing such series. As one increases
the order of the Maclaurin polynomial, the approximation to sin x becomes increasingly
better
1st - order Maclaurin polynomial : x cos(0) = x
x2
2 nd - order Maclaurin polynomial : x cos(0) sin(0) = x
22!
x
x3
x3
3 rd - order Maclaurin polynomial : x cos(0) sin(0) cos(0) = x
2!2
3!3
3!
4
x
x
x
x3
4 th - order Maclaurin polynomial : x cos(0) sin(0) cos(0) +
sin(0) = x
2!
3!
4!
3!
x2
x3
x4
x5
x3 x5
th
+
5 - order Maclaurin polynomial : x cos(0) sin(0) cos(0) + sin(0) + cos(0) = x
2!
3!
4!
5!
3! 5!
Notice in the graph below that as the degree of polynomial increases, so does its
nearness to sin x in the vicinity of the origin:
62
The dark, shaded curve is y = sin x , the solid line is y = x, the dotted curve is
x3
x3 x5
+ . It should be clear that as the degree
, and the solid curve is y = x
3!
3! 5!
of the polynomial increases, so does its agreement with the function that it approximates
in the vicinity of where the Maclaurin series is centered (i.e. x = 0 ).
When solving problems concerning Maclaurin series, it is helpful to make a table
with one column for n, one column for f ( n ) ( x), and one column for f ( n ) (0) to ease the
burden of constructing the Maclaurin polynomial. For instance, if one were to use a 3rdorder Maclaurin polynomial to approximate the value of sin x at 0.01, he or she would
create the following table:
n
f ( n ) ( x)
f ( n ) (0)
sin x
0
0
cos x
1
1
sin x
2
0
cos
x
3
-1
y = x
Using the table as a guide, one can easily generate the Maclaurin polynomial:
n =0
(0.01) 3
= 0.0099998333.
3!
Ex.) Using a 3rd-order Maclaurin series approximate the value of tan(0.2) .
n
0
1
2
f ( n ) ( x)
tan x
sec 2 x
2(sec 2 x)(tan x)
n =0
f ( n ) (0)
0
1
0
4
2(0.2) 3
= 0.1946666667.
3
As a graphical supplement, observe the proximity to tan x of the polynomial in
the vicinity of x = 0 :
tan(0.2) (0.2)
63
The solid curve is tan x , and the dotted curve is the polynomial.
Taylor Series: When Colin Maclaurin published the series that are now named after him,
he knew very well of the toils of an English mathematician named Brook Taylor. In fact,
Maclaurin series are merely special cases of Taylor series; while the former are centered
at the origin, the latter can be centered anywhere. Taylor series are thus used when
approximating values of functions that are not in the vicinity of x = 0. For instance,
while Maclaurin series are valid for x-values such as 0.01 and 0.2, they are not valid for a
value such as x = 2. In such cases, one approximates the complicated function in
question with a Taylor polynomial, which is generated by the following rule:
( x c ) n f ( n ) (c )
( x c) 2 f (c) ( x c) 3 f (c)
= f (c ) + ( x c ) f (c ) +
+
+ L.
n!
2!
3!
n =0
Note that this is essentially the same formula as that of the Maclaurin series,
except that this formula is centered on a non-zero value c.
Ex.) Approximate the value of ln(1.6) using a 5th-degree Taylor polynomial
centered at x = 2.
Just as with the setting up of Maclaurin polynomials, a table greatly facilitates the
creation of Taylor polynomials:
n
0
1
2
3
4
f ( n ) ( x)
ln x
1/ x
1/ x 2
2 / x3
6 / x4
f ( n ) (2)
ln(2)
1/ 2
1/ 4
1/ 4
3/8
24 / x 5
3/ 4
( x 2) 0 ln 2 ( x 2)(1 / 2) ( x 2) 2 (1 / 4) ( x 2) 3 (1 / 4) ( x 2) 4 (3 / 8) ( x 2) 5 (3 / 4)
+
+
+
+
+
0!
1!
2!
3!
4!
5!
2
3
4
5
( x 2) ( x 2)
( x 2)
( x 2)
( x 2)
= ln 2 +
+
.
2
8
24
64
160
64
+
= 0.47001651.
2
8
24
64
160
Notice the close proximity of y = ln x and the Taylor polynomial in the vicinity of
ln(1.6) ln(2) +
x = 2:
The solid curve is y = ln x and the dotted curve is the Taylor polynomial.
In the previous chapter, the term truncation was introduced in the context
of splitting an infinite series into a certain number of partial sums while leaving
the rest of the series unaccounted for. In performing such an action, one can
determine the error associated with the chosen partial sum. Truncation and the
error associated with it can be interpreted similarly for Taylor series. However,
one must understand truncation in the context of a very important theorem. In
one of Taylors most influential papers, Methodus incrementorum directa et
inversa (1715), he describes what is now known as Taylors theorem. This
theorem concerns the relationship between the chosen Taylor polynomial, denoted
as Pn (x), and the rest of the infinite series, the remainder, denoted as Rn (x).
Taylors theorem states that if a function f (x) has n derivatives on the closed
interval a x b and its (n + 1) st derivative exists on the open interval a < x < b ,
there exists some number (the Greek letter xi-pronounced zye) in the closed
f ( n +1) ( )( x c) n +1
and
(n + 1)!
f ( x) = Pn ( x) + Rn ( x). Now for the English version: Taylors theorem states that
a function f (x) that may be defined as a power series consists of the chosen
polynomial, Pn (x), and the remainder Rn (x), which makes intuitive sense; there is
a component of the series that one is using and a component that one is not using.
The remainder term Rn (x) is nonetheless quite useful because it allows one to
determine the error associated with the polynomial approximation. It is important
to understand the significance of the terms in the formula for the
f ( n +1) ( )( x c) n +1
remainder: Rn ( x) =
. The term f ( n +1) refers to the (n + 1) st
(n + 1)!
derivative of the function that one is approximating. For instance, if one were
using a 4th-order Taylor polynomial, he or she would evaluate the 5th derivative of
that function for the remainder formula. Note that this derivative is evaluated at a
65
number , which some books may refer to as z. This number occurs between the
center of the series, c, and the value that one is evaluating, x. In order to
determine the maximum possible error associated with a polynomial
approximation, it is necessary to choose a value of between x and c such that it
makes the error, Rn (x), as large as possible. The formula
f ( n +1) ( )( x c) n +1
represents the Lagrange form of the remainder,
(n + 1)!
named after Joseph Louis Lagrange, an 18th-century Italian-French
mathematician. There is also another form known as the Cauchy form of the
remainder, named after the 19th-century French mathematician Augustin Louis
Cauchy, which involves integration, but will not be discussed here.
Notice the similarity between this error analysis and that which was
introduced in the previous chapter concerning alternating series. Both involve the
use of the (n + 1) st term to determine the truncation error. In the context of Taylor
series, the truncation error, the absolute value of the remainder, must be less than
or equal to the Lagrange term. The following two examples should clarify the use
of the Lagrange form of the remainder to estimate.
Rn ( x ) =
n
0
1
2
3
4
ex
ex
ex
ex
ex
f ( n ) (1)
E
E
E
E
E
e( x 1) 0 e( x 1)1 e( x 1) 2 e( x 1) 3 e( x 1) 4
+
+
+
+
0!
1!
2!
3!
4!
2
3
4
e( x 1)
e( x 1)
e( x 1)
= e + e( x 1) +
+
+
2!
3!
4!
e(0.8 1) 2 e(0.8 1) 3 e(0.8 1) 4
e e
e
e
e
e + e(0.8 1) +
+
+
=e +
+
=
2!
3!
4!
5 50 750 15000
2.225547942.
( 0.8 )
( )(0.8 1) 5
. What value of
5!
must be chosen? In order to make the error as large as it can possibly be, one chooses
the largest value of f(x) on the closed interval 0.8 x 1 , which is 1. Thus,
f ( 5) (1)(0.2) 5
e ( 0.2) 5
Max error Rn (0.8) =
=
7.24875154 10 6.
5!
5!
Taylor polynomial was used, (n+1) = 5, and Rn (0.8) =
(5)
In another type of problem, one must determine the order of the Taylor
polynomial to be used based upon the maximum acceptable error.
Ex.) What minimum order Taylor polynomial centered at x = / 2 must be used
to approximate the value of cos( / 5) with an error no greater than 0.0003?
In this problem, c = / 2 and x = / 5. Since n is not known, neither is
st
the (n + 1) derivative known. Thus, the appropriate value of must be chosen
somewhat more hypothetically than in the previous problem. Another complexity is that
differentiation of cos x makes for an oscillation between positive and negative values of
cos x and sin x. In any event, what is the greatest value that a pure (i.e. without
multiplying, dividing, adding, or subtracting constants) sine or cosine function may have?
The answer 1. Which function (i.e. cos x, cos x, sin x, or sin x) has a value of 1
somewhere on the closed interval / 5 x / 2 ? The only one is sin x , and this occurs
at x = / 2. Thus, the (n + 1) st derivative is sin x, and the value for is / 2. However,
the value of n is still not known, but this can be easily found from the formula for the
f ( n +1) ( )( x c) n +1
Lagrange form of the remainder: Rn ( x) =
:
(n + 1)!
Max error Rn ( x) =
f ( n +1) ( )( x c) n +1
(n + 1)!
sin( / 2)( / 5 / 2) n +1
(0.0003) Rn ( / 5) =
(n + 1)!
(3 / 10) n +1
. This would be relatively simple to solve if it were
(n + 1)!
not for the factorial, so one must use trial and error. By testing values of n, one finds that
(3 / 10) ( 6+1)
(3 / 10) (5+1)
n must equal 6, because
= 0.0001311 and
= 0.0009734 .
(6 + 1)!
(5 + 1)!
Thus, one must use a 6th-order Taylor polynomial.
(0.0003)
An important technicality must be introduced at this point. Thus far, it has been
assumed that the Taylor polynomials used all converge within the area in which the
approximation is taking place. For functions such as cos x, sin x, e x , ln x , and other
common functions, this is true. In general, however, convergence is not necessarily
67
guaranteed. However, one can determine if a Taylor series converges in the neighborhood
in which one is approximating by using the formula for the Lagrange form of the
remainder. As the order of the Taylor polynomial increases, the remainder should
become smaller. This should make sense when considering that f ( x) = Pn ( x) + Rn ( x); as
n increases, the term Pn (x) will become larger, which means that the term Rn (x) must
become smaller. In fact, for a convergent Taylor series, when n = , Rn ( x) = 0. This
should also make sense; if an infinite-order Taylor polynomial is used to approximate,
there should be no remainder, and no error! Thus, if lim Rn ( x) = 0, the Taylor series is
n
convergent in the vicinity of c, and perhaps everywhere. This intricacy, though important
theoretically, is not an aspect of the AP Calculus BC curriculum.
in
which
sin
,
is the angular acceleration of the
l
dt 2
dt 2
pendulum, g is the acceleration due to gravity, l is the length of the pendulum, and is
68
the angular position of the pendulum. Unfortunately, this equation does not represent
simple harmonic motion, which would state that the second derivative of position (i.e.
acceleration) is equal to a negative constant times the position, which would have the
d 2
g
form of
= . It is the case, however, that the swinging of a pendulum can be
2
l
dt
approximated as simple harmonic motion for small angles. Why is this the case? For
small angles, sin . Thus, as long as the pendulum does not stray too far from the
d 2
g
vertical, it can be described by the equation
= . To understand this
2
l
dt
graphically, consider the fact that the line y = x approximates the value of y = sin x very
well be in the neighborhood of the origin:
d 2
g
This is an important because, relatively speaking, the expression
= is a
2
l
dt
second-order differential equation that can be easily solved while the expression
d 2
g
= sin is basically intractable. In essence, one is performing a polynomial
2
l
dt
approximation to simplify the differential equation at hand. There are many other
situations in which polynomial approximations are necessary in the physical sciences, as
will shown in the next two subsections.
Modeling Molecular Vibrations: All objects in the universe exert forces on one
another. This is true for matter at the submicroscopic level as well. Molecules are
composed of atoms participating in chemical bonds, and these atoms exert both repulsive
and attractive forces upon other nearby atoms. Experiments in the field of quantum
chemistry by the English mathematician John Lennard-Jones in 1931 have led to a
function that conveys how the potential energy of the atoms in a noble gas (e.g. helium
and argon) varies with increasing distance between these atoms. This function, known as
12 6
the Lennard-Jones potential, is given as V (r ) = 4 , where V (r ) is the
r
r
potential energy as a function of the distance between the atoms in a molecule and and
are both constants. Below is a plot of the Lennard-Jones potential, with distance r on
the x-axis and potential V on the y-axis:
69
The Lennard-Jones potential, as one can see from the unwieldy formula, is rather
difficult to work with in computational studies. Thus, one chooses to simplify it as a
Taylor polynomial that approximates the potential energy values in the vicinity of the
distance associated with the lowest potential energy value. The Taylor polynomial
expansion of the Lennard-Jones potential about is:
n =0
(r ) n V ( n ) ( )
(r ) 2
(r ) 3
= V ( ) + V ( )(r ) + V ( )
+ V ( )
+L
n!
2!
3!
Noting that is the minimum x-value, this is the value at which the first
derivative of the Lennard-Jones potential point should be zero. This allows one to find
in terms of :
6 6
V ( ) = 4
7
( )
12 12
13
( )
= 0 6 6 13 = 12 12 7 6 13 = 2 12 7
6
7
=
= 26 .
12
13
2
Thus, the first two terms of the Taylor series are
12 6
V ( ) = 4 6 6 = and zero, since the first derivative at was
2
2
42 6 36 3 4
(2 6 ) 8 = 2 .
36 (r ) 2
.
Thus, the Taylor polynomial is V (r ) +
3
2 2
Observe how closely, even at just two terms, this polynomial approximates the
Lennard-Jones potential in the vicinity of :
70
The thicker curve is the Lennard-Jones potential and the thinner curve is the
second-order Taylor polynomial.
Given the relatively cumbersome calculations involved here, one might ask,
How did finding a Taylor polynomial help us at all? The answer lies in a concept
known as computational cost. While the initial calculations might have seemed
substantial, the resulting polynomial eases subsequent calculations, either by hand or by
computer, to such a degree, that the net payoff is considerable.
The Virial Equation: To the readers who have studied chemistry and physics, the
equation PV = nRT must look familiar. This is known as the ideal gas law, in which P is
pressure, V is volume, n is the number of moles of the gas (a measure of the number of
particles), R is a constant, and T is temperature. As the name of this equation suggests, it
is an idealization that only works in rather limited circumstances. Many great scientists
have modified the ideal gas law to describe the behavior of real gases, but even these
equations are often insufficient. One need not approach the problem with knowledge of
the explicit algebraic structure of the mathematical model, however. In fact, one of the
powerful attributes of Taylor series is that they can function as mathematical models even
when one does not know anything about the actual structure of the equation at hand. In
the case of a real gas, the pressure (P) of the gas in question can be described as a
function of the concentration of the gas particles (M) and temperature (T): P = f ( M , T ).
Notice that this equation is not very specific at all. In fact, it merely states that there is
some function that expresses the relationship between the three variables. This is of little
concern. Assuming that the function is continuous, one may expand it as a Taylor series
1 2 f (M , T ) 2
f ( M , T )
as follows: P = f ( M , T ) +
M+
M +L
2! M 2
M
Do not be put off by the new notation; the symbol means that one is taking what is
known as a partial derivative, because the function in question is one of three variables,
not just two. A (much) more extensive discussion of partial derivatives is reserved for
multivariable calculus. In any event, it is important to note important aspects of the
Taylor polynomial above. First of all, the partial derivatives are merely coefficients of
the variable term M. Secondly, one can evaluate this polynomial at a gas concentration
of zero, which means that M = 0, technically turning the Taylor polynomial into a
Maclaurin polynomial. When the gas concentration is zero, the gas becomes ideal,
n
meaning that the first term, f ( M , T ) , is RT , which can be written as MRT , because M
V
f ( M , T )
2 f (M , T )
just refers to the moles of gas over volume. Since the terms
and
M
M 2
71
are merely coefficients, one can simplify the polynomial by dividing both sides by MRT
and giving the partial derivatives simpler names of C and B:
P
= 1 + BM + CM 2 + L.
MRT
The equation above is known as the virial equation, and B and C are virial coefficients,
which can be determined experimentally or theoretically. One need not know the science
behind this equation to appreciate its significance; almost nothing was known before the
polynomial expansion, but now this equation yields very useful information about the gas
in question, more useful, in fact, than could any of the other algebraically explicit
equations!
In these last two subsections of this chapter, the application of Taylor series to
problems that often arise in mathematics will be discussed.
Taylor Series and Integration: It is often the case that one will encounter integrals that
simply cannot be expressed in terms of familiar functions. Not only are these integrals
difficult to evaluate, but it is impossible to evaluate them in terms of algebraic or
transcendental functions. Consider the following definite integral:
is actually no familiar function whose derivative is sin( x ). (If the reader really must
know, this kind of integral is called a Fresnel integral, which appears often in wave
optics). However, integration of polynomials is the easiest integration there is! The
definite integral above can be approximated as the definite integral of a Maclaurin
polynomial (because the values over which one is integrating are close to the origin).
n =0
x3 x5 x7
(1) n x 2 n +1
= x
+
+ L.
(2n + 1)!
3! 5! 7!
One can simply substitute x for x , allowing one to approximate the definite integral as
follows:
1
1
x 6 x10 x14
sin( x 2 )dx = x 2
+
+ Ldx .
3!
5!
7!
0
0
Taylor Series and Differential Equations: The very broad topic of differential equations
was introduced in chapter 2. There are countless analytical and numerical techniques that
are used to solve differential equations. In chapter 2, the method of separation of
variables and Eulers method was discussed. This subsection introduces another
technique, known as the method of power series, which is quite helpful in solving
second-order differential equations. This method is best explained through a simple
dy
example. Consider the differential equation
= ky. From chapter 2, it was determined
dx
that the solution is an exponential function. While this differential equation can be easily
solved by the method of separation of variables, it is used as an example here to provide a
72
dy
= ky
dx
by the method of power series. To do this, one would first perform a polynomial
expansion upon y as follows: y = a 0 + a1 x + a 2 x 2 + a3 x 3 + K , in which the a coefficients
are merely the constant terms of the power series. With an initial condition, one can
determine the value of a0 . In this case, y (0) = 1, so y = 1 + a1 x + a 2 x 2 + a3 x 3 + K . After
this step, one determines the polynomial expansion for y , which is simply the derivative
of the polynomial expression already found: y = a1 + 2a 2 x + 3a3 x 2 + L . One can
actually equate these two polynomial expansions in the context of the given differential
dy
equation:
= kx a1 + 2a 2 x + 3a3 x 2 + L = k a 0 + a1 x + a 2 x 2 + a3 x 3 + K
dx
For the sake of simplicity, assume that k = 1. Now one can equate the coefficients on
each side of the equation:
a1 = a0 = 1
1
1
1
2a 2 = a1 a 2 = a1 = (1) =
2
2
2
1
11
1
3a3 = a 2 a3 = a 2 = =
3
3 2 3 2
M
M
M
M
N
1
1
na n = a n 1 a n = a n 1 =
n
n!
2
3
n
x
x
x
Thus, y = 1 + x +
+
+L
= ex.
2! 3!
n!
With these steps in mind, one can solve the following second-order differential
d2y
= ky. This second-order differential equation is important because it is
equation:
dx 2
the kind that describes simple harmonic motion, described earlier in this section. One
must determine polynomial expansions for y, y , and y :
y = a 0 + a1 x + a 2 x 2 + a3 x 3 + L . Given the initial condition that y (0) = 1,
simplified context for the method of power series. Suppose one were to solve
y = 1 + a1 x + a 2 x 2 + a3 x 3 + L.
y = a1 + 2a 2 x + 3a3 x 2 + L.
y = 2a 2 + 6a3 x + L.
These polynomials can be equated to each other in the context of the differential
d2y
= ky (2a 2 + 6a3 x + L) = k (1 + a1 x + a 2 x 2 + a3 x 3 + L).
equation:
2
dx
For the purpose of simplification, assume that k = 1 :
( 2 a 2 + 6 a 3 x + L) = 1 a1 x a 2 x 2 a 3 x 3 L .
Equating the coefficients:
1
2 a 2 = 1 a 2 =
2
a
6a3 = a1 a3 = 1
6
( 1 / 2) 1
12 a 4 = a 2 a 4 =
=
12
24
a3
a1
20 a5 = a3 a5 =
=
20 120
73
a
1
a
1
y = 1 + a1x x 2 1 x 3 + x 4 + 1 x 5 + L
120
24
6
2
Note that since there is no a1 on the left side of the original equality, yet one on
the right side, a1 must equal zero. Thus,
1
1
y = 1 x2 + x4 +L.
2!
4!
One recognizes this as the power series expansion for y = cos x. In fact, this is what
one would expect from an equation that models simple harmonic motion. Imagine a
pendulum swinging back and forth at small angles. If one were to record its position
versus time, a resulting plot of the data would yield a cosine or sine function.
As always, the method of power series does have certain limitations. This method
can only be used when the power series used are centered at the origin. Also, it can only
be used for what are known as homogeneous differential equations, in which the only
terms are the function in question and a certain number of derivatives, all multiplied by
d2y
constants. For instance, the differential equation
= ky is homogeneous because
dx 2
the only terms are the function in question and its second derivative multiplied by
d2y
= ky + 1 , however, is not homogeneous
constants. The differential equation
dx 2
because it includes a term (1) that is not the function in question or its derivatives
multiplied by constants. A more extensive coverage of this material is the stuff of a
course in ordinary differential equations, but one would ordinarily have to take Calculus
III and Linear Algebra before it.
Concluding Remarks
This chapter introduced the theory and applications of a very important
mathematical tool power series. These expressions are similar to the infinite series
covered in the previous chapter, only power series contain variable terms. In addition,
when determining the convergence of a power series, one must specify where it
converges. Two important types of power series, Maclaurin and Taylor series, were
discussed in the context of polynomial approximations. The technique of polynomial
approximations to difficult functions was also covered in greater depth through a
discussion of their indispensable applications to problems in science and mathematics.
Key Terms:
power series
radius of convergence
interval of convergence
term-by-term differentiation
term-by-term integration
Maclaurin series
Maclaurin polynomial
Taylor series
computational cost
Taylor polynomial
virial equation
Taylors theorem
method of power series
remainder
Lagrange form of the remainder
simple harmonic motion
Lennard-Jones potential
74
The infinitesimal line segment dl is the hypotenuse of the triangle formed by the
perpendicular line segments dy and dx. Using the Pythagorean theorem:
dl 2 = dy 2 + dx 2 dl = dy 2 + dx 2 . Casting this equation in a form that contains the
75
dy 2
dy 2
dx
1 + 2 = 1 + dx . Many
dx
dx
dy
recognizable derivative
, dl =
dx
dy 2
1 + dx .
a
dx
The above is the formula for arc length (l). Oftentimes, this integral is far too
complicated to evaluate analytically by conventional methods, so the graphing calculator
is often used for arc length problems. This formula can be used to prove that the
circumference of a circle is indeed 2 r . Consider the formula for a semicircle with a
l=
2x
dy
=
=
dx 2 r 2 x 2
One then uses the formula for arc length:
l=
Use
x
1 +
r 2 x 2
du
a2 u2
= sin 1
dx =
x2
1 + 2
2
r x
dx =
x
r x2
2
r 2 x2
x2
2
+
2
r 2 x2
r x
dx =
r2
2
2
r x
dx
u
+ C. Let u = x, du = dx, a = r :
a
r
r
r
r2
r2
r
x
r
r
2 2 dx =
= rsin1 rsin1 = rsin1 (r) rsin1 (r)
dx =
dx = rsin1
2
2
2
2
r r
r r x
r r x
r r x
r
r
= r . This is the length of the semicircle, so the length of the whole circle must
be 2 r. It appears as though the ancient mathematicians were right.
r
76
l=
of these, Y1 , represents 1 + ln 2 x cos 2 x . Select Y1, which should now appear in the Y2
ln x cos x sin x
." After scrolling
slot on the Y= screen. In the Y2 slot, type Y1 + 2
x
down to the Y3 slot, go to the VARS menu as before and select Y2, which should now be
sin 2 x
. Finally,
displayed on the Y3 slot on the Y= screen. In this slot, type Y2 +
x2
scroll down to the Y4 slot and press the 2ND and x2 buttons, which should put a radical
symbol in the Y4 slot. Go to the VARS menu, selecting Y3. This should put the Y3
within the radical symbol in the Y4 slot. The function Y4 actually represents
ln x cos x sin x sin 2 x
1 + ln x cos x + 2
. Overall, the Y= screen should look like
+
x
x2
this:
2
Use the method explained above whenever a structurally difficult problem must
be solved graphically.
Now, one can evaluate the integral of Y4 from x = 1 to x = 2 :
2
77
l=
dx 2
1 + dy ( Where c and d are the y - values associated with a and b).
dy
l=
since the calculator does not know the difference between y(x) and x(y). Solving with the
graphing calculator, one gets l = 6.3083997.
Arc Length With Corners or Cusps: Recall that a function is not integrable if there are
corners ( ) or cusps ( { ) along the interval of integration. Thus, if the derivative of
the function whose arc length one is trying to find has any of these points, one must
work around them by making a piecewise function. Consider, for instance, the
5
2
5 3
3
function y = x . The derivative of this function is x , which has has a cusp at x = 0 :
3
Thus, the original function must be defined as a piecewise function about the
53
origin: y = x 5 , x 0 One must then integrate the first piece from the lower limit to 0
x 3 , x 0.
and then integrate the second piece from 0 to the upper limit.
5
Ex.) Determine the length of the curve y = x 3 on the interval [-2, 2].
One integrates each piece separately and then adds each numerical result
to yield the total arc length:
78
l=
5 23
1 + x dx +
3
5 23
1 + x dx = 3.8493405 + 3.8493405 = 7.698681.
3
This is exactly the path that one would expect from throwing a ball in the air a
parabola. But why is nearly half of the parabola missing in the graph? This concerns a
concept known as the parameter interval, which is the domain of the parameter. The
calculator has, in the case of the curve above, automatically made this interval
0 t 2 . Thus, the initial point of the parabola occurs at (x(0), y (0) ) and the
terminal point of the parabola occurs at (x(2 ), y (2 ) ). In this case, one does wish for
the initial point to occur when t = 0 because t is the time parameter. In some cases,
however, it may be necessary to define a specific parameter interval. One can change the
79
parameter interval on the graphing calculator by going to the WINDOW menu while in
parametric mode.
Converting Between Parametric and Rectangular Form: Sometimes, it is
desirable to convert from parametric form back to rectangular (Cartesian) form. While
this can often be done algebraically, there are many cases in which this is very difficult.
The parametric equations already discussed are an easy example. Given x(t ) = 1 + t and
y (t ) = 1 + t t 2 , one can find y as a function of x:
x = 1+ t t = x 1
y = 1 + t t 2 y = 1 + ( x 1) ( x 1) 2 = 1 + ( x 1) ( x 2 2 x + 1) = x 2 + 3 x 1.
It is almost certainly the case that one will not be expected to derive parametric
form from rectangular form, as this often requires either a lot of mathematical theory, or
a lot of mathematical experiment. There are some certain cases in which this is not that
difficult, but it is not expected on the AP Calculus BC exam. On the exam, one will be
given parametric equations and asked to perform an operation on them or apply them in a
certain way. Not only will one not be asked to convert from rectangular form to
parametric form, but it is also very rare that one will be asked to convert from parametric
form to rectangular form. Below are some examples of common curves seen in
parametric form:
The Circle: This is the only curve for which a derivation of the conversion
between Cartesian and parametric coordinates will be shown. Consider a particle moving
along the circle below:
The parameter t represents an angle in this case. One can define the particles xand y-positions in terms of this parameter through the use of right triangle trigonometry.
If a is the hypotenuse:
x = a cos t and y = a sin t.
Thus a circle has the parameterized form of x(t ) = a cos t , y (t ) = a sin t , where a is
the radius of the circle. For instance, the parametric equations x(t ) = cos t , y (t ) = sin t
represent the unit circle.
The Ellipse: An ellipse has the parameterized form x = a cos t , y (t ) = b sin t. If
a > b, then a is referred to as the major axis and b is referred to as the minor axis. If
a < b, the reverse is true. For instance, the ellipse x(t ) = 2 cos t , y (t ) = sin t is shown
below:
80
The major axis is 2 and the minor axis is 1. Note that a circle is merely an ellipse
with a = b.
The Astroid: At this point, more esoteric parametric curves will be considered, but
they all may very well appear on the AP Calculus BC exam. An astroid (not asteroid)
is a parametric curve generated by tracing the path of a point at the edge of a circle
rolling inside of another circle. The general parametric form of an astroid is
x(t ) = a cos 3 t , y (t ) = a sin 3 t , in which a is the radius. For instance, the astroid
x(t ) = cos 3 t , y (t ) = sin 3 t is shown below:
The Cycloid: A cycloid is a parametric curve defined by the path of the point on
the edge of a circle that rolls along a straight line. It has the general parametric form
x(t ) = a(t sin t ), y (t ) = a (1 cos t ). Note that in the cycloid, unlike in other situations,
the x-coordinate is associated with the sine function and the y-coordinate with the cosine
function. The cycloid played a prominent role in physics during the time of Newton
because it was used to find solutions to the brachistochrone and tautochrone problems.
The former concerned the curved path between two points in which an object covered the
distance of the curve in the least amount of time, and the latter concerned the curious
observation that an object in curved space, such as a bowl, will take the same amount of
time to reach the bottom regardless of its starting point on the curve. As an example,
consider the cycloid x(t ) = 3(t sin t ), y (t ) = 3(1 cos t ) shown below:
81
The Witch of Agnesi: The Witch of Agnesi is an example both of parametric form and of
linguistic error. This curve was discovered by the 18th-century Italian mathematician
Maria Gaetana Agnesi. It is described by the parametric equation
x(t ) = 2 cot t , y (t ) = 2 sin 2 t , which shown graphically below:
82
dy
2
dy
dy dt 2 sin t
=
=
= tan t
5 cos t
5
dx
dx dx
dt
t=
2
3
2 2
= tan
= 0.692820323.
5 3
The slope of the tangent line to a parametric curve has an important interpretation
in physics. In the context of particle motion, the derivative of a parametric equation
represents the direction in which the particle is moving. For instance, a derivative of zero
implies that the particle is moving horizontally and a derivative of infinity implies that
the particle is moving directly up or down.
Ex.) A particle is moving in two dimensions as described by the parametric
equations x(t ) = sin t , y (t ) = 2t 2 5. At t = 3.5, what is the slope of the tangent line and
in which direction is the particle moving?
dy
dy
dy dt
=
= 4t sec t
= 4(3.5) sec(3.5) = 14.94997066. The particle
dx t =3.5
dx dx
dt
is moving upwards to the left.
To see why this is the case observe the parametric curve below, which represents
the particles path:
The black dot represents the particle, which is moving upwards to the left.
The evaluation of the second derivative of a parametric equation is slightly more
difficult than finding the first derivative. The second derivative ( d 2 y / dx 2 ) is not merely
the quotient of the second derivatives of the x- and y- equations with respect to t. Instead,
one must apply the chain rule in a somewhat less obvious way:
dy
d
dx
2
d y d dy d dy dt
= =
= dt . Here is an algebraically difficult
2
dx
dx dx dt dx dx
dx
dt
problem:
83
dy
:
dx
dy
2 sin t cos t
dy dt 4 sin t cos t
=
=
=
= 2 sin 3 t cos t.
2
1
dx dx
2 csc t
dt
sin 2 t
Next, determine
d dy
:
dt dx
d dy
3
3
= 2 (sin t )( sin t ) + (cos t )(6 sin t cos t )
dt dx
Finally, divide this expression by
dx
:
dt
dy
d
dx
3
3
3
3
2
dt = 2 (sin t )( sin t ) + (cos t )(6 sin t cos t ) = (sin t )( sin t ) + (cos t )(6 sin t cos t ) = d y .
dx
dx 2
csc 2 t
2 csc 2 t
dt
] [
l=
84
l=
dy 2
1 + dx =
dx
b
a
t2
t1
b
a
dy / dt 2
1 +
dx =
dx / dt
(dx / dt ) 2 + (dy / dt ) 2 dx
dt =
dt
(dx / dt )
2
t2
t2
t1
b
a
(dx / dt ) 2 + (dy / dt ) 2 dx
dt
(dx / dt ) 2
dt
2
dx
dy
+ dt
dt
dt
dx
dy
l =
+ dt , where t1 is the value of t associated with a and t 2 is the
dt
dt
t1
value of t associated with b. Note that in the context of particle motion, this
parameterized formula refers to the total distance traveled by the particle.
4
dx
dy
( 6 sin t )2 + (3 cos t )2 dt.
+
dt
=
dt
dt
t1
0
Evaluating with the graphing calculator, which should be done in FUNCTION
l=
t2
mode,
l=
36 sin 2 t + 9 cos 2 t dt = 17.783441. Note that the integral could have been
easily evaluated analytically, but one would had to have used the calculator eventually
when plugging in a value of 4 for t.
Area Under Parametric Curves: It is hardly ever the case that one is asked to
determine the area enclosed by a parametric curve, and certainly not on the AP Calculus
BC exam. However, since this information is so scant, I thought it appropriate to include
it in this book. The derivation is really quite simple. Recall that the area under a curve
x2
x1
t2
t1
85
Area =
1 + cos 2
. Remember that the common
2
trigonometric identities must be memorized for the AP Calculus BC exam. They can be
found in appendix A of this book.
2
2
2
27
9 9 cos 2t
9
18 cos t + 9 dt =
(9 cos 2 t 18 cos t + 9)dt =
+
cos 2t 18 cos t + dt
2
2
0
0 2
0 2
Use the trigonometric identity cos 2 =
27
1
= sin t 18 sin t t = 27 .
2
2 0
dx
dy
2 y + dt
dt
dt
a
dx
dy
2 x + dt .
dt
dt
a
Notice that the radical expression is the arc length, which is multiplied by another
linear term to yield the surface area.
(When revolved about the x-axis): Surface Area =
dx
dy
2x + dt . What are the limits of integration? Starting at t = 0, one
dt
dt
a
reaches this point again at t = 2 . Thus, these are the two limits of integration.
86
Surface Area =
2
0
Using the graphing calculator in FUNCTION mode, one can easily evaluate this integral:
Surface Area = 157.91367. This is the surface area of one doughnut.
Multiplying by 25 yields the total area of icing that the baker needs:
Total Surface Area = 25157.91367 = 3947.84175 square units of icing.
Introduction to Polar Curves
Unlike Cartesian and parametric coordinates systems, there is a system
that defines points in space not by x- and y- coordinates, but by lines and angles. This
coordinate system is associated with polar equations. A point in space in a polar
coordinate system is defined by a directed distance (r) of certain length and the angle
(in radians) at which this distance is displaced from a reference ray, known as the initial
ray. Thus a point P has the coordinative notation of P(r , ) . This is supported by the
diagram below:
87
r = x2 + y2
y
= tan 1 , x 0.
x
Ex.) Express the equation xy = 4 in polar form.
Since x = r cos and y = r sin , (r cos )(r sin ) = 4
r 2 cos sin = 4.
Polar Curves: Polar curves are defined as the directed distance r as a function of (i.e.
r ( ) . Thus, as r and vary, they define a curve in polar space. For instance, consider
the simplest polar curve, a circle. A circle has a constant r with an angle that varies.
Thus, the polar equation for a circle is simply r ( ) = constant, where the constant is the
radius. To graph in polar form on the TI-84, go to MODE and choose POL. In the Y=
screen, one should now see a list of r =. To graph a circle with a center at the origin
and a radius of 3, simply input 3 for r1. The result is:
88
It will almost certainly not be required to produce a polar curve on polar graph paper
on the AP Calculus BC exam. Even if one is asked to reproduce a polar curve on paper,
it will most likely be in the context of regular graph paper.
Obviously, polar curves become more elaborate once one allows r to vary along with
. Below is a discussion of the various sorts of polar curves that one is most likely to
encounter:
Spirals: In polar coordinates, a spiral can be viewed as a curve that turns around a
point of reference while getting closer or farther away from it, depending upon the
equation representing the spiral. Ancient mathematicians were fascinated by these
geometric phenomena because they are often manifested in nature. Many of the spirals
discussed here were first used to model these natural systems.
The Archimedean spiral has the general form r ( ) = a + b , where a and b are
constants. For instance, the Archimedean spiral r ( ) = 1 + 2 is shown below:
89
Note that the calculator automatically makes the interval of the angle [0, 2 ] . To
view more of the spiral, one could make max larger in the WINDOW option. For
example, allowing max to equal 4 yields a fuller spiral:
A key feature of the Archimedean spiral is that the distance between successive turns
is always the same. This is not the case for the logarithmic spiral, which has the general
form r ( ) = ae b , where a and b are constants. For instance, consider the logarithmic
spiral r ( ) = e
For this graph max was made 8 . The turns near the origin are so small compared to
the noticeable turns that they only appear as a dark spot. The primary observation to take
away from this is that at each subsequent turn, the distance of that turn relative to the one
before it increases exponentially. Logarithmic spirals are extremely prevalent in nature.
They are seen, for instance, as mollusc shells, spiral galaxies, hurricanes, and the cochlea
(an organ of the inner ear) of mammals. There are many other spirals that can be
expressed in polar form, but they are slightly more esoteric and will not be covered here.
The Limaons: A limaon is a polar curve with the general form r ( ) = a b sin or
r ( ) = a b cos , where a and b are constants. The term comes from the Latin limex,
which means snail, referring to the general shape of these polar curves. The limaon
was first described by the German Renaissance artist Albrecht Drer in his
Underweysung der Messung (Instructions in Measurement) (1525). The limaon family
is generally composed of four types of curves, which are shown below:
90
convex limaon
Conditions: 2a b
cardioid
Conditions: a = b
Note that memorization of these terms is not required for the AP Calculus BC exam,
but it is helpful to know them nonetheless. The cardioid is so well-studied among the
limaons that it has its own name! It is a very common polar curve in which a = b.
Geometrically, a cardioid looks similar to a heart shape, hence its name. The cardioid is
generated by tracing the point on the edge of a circle as it rolls along a circle of the same
radius. For instance, consider the cardioid represented by r ( ) = 2 + 1.5 sin :
In microphone science, the cardioid is the most common shape for the sounds picked up
by a unidirectional microphone. A unidirectional microphone is sensitive only to sounds
coming from a specific direction, which is helpful when one does not want to pick up
extraneous noise.
Besides the cardioid, the AP Calculus BC exam often includes the limaon with a
loop, in which a < b . For instance, the limaon represented by r ( ) = 1 + 2 cos is
shown below:
As will be discussed in upcoming section, the AP exam often likes to have students
analyze the loop in some way.
The Lemniscate: A lemniscate is geometrically a figure-eight. It was first described by
Jakob Bernoulli (a member of that famous mathematical dynasty) in 1694. The name is
derived from the Latin lemniscus, which means ribbon. It is also sometimes used to
91
describe the symbol for infinity ( ) . It is described by the general polar equation
r 2 = a 2 cos 2 , where a is a constant. Unfortunately, the lemniscates cannot be graphed
in this form on the TI-84; one must put the equation in as r ( ) = a cos 2 . Thus, in
one slot on the Y= screen, one would input the positive value, and in another slot the
negative value. For instance, the lemniscate represented by r 2 = 25 cos 2 would be
written as r ( ) = 5 cos 2 , which is shown below:
Unfortunately, this is the best the calculator can do. The space where the calculator
thinks no curve exists is due to taking the square root of a negative number at those values, yielding imaginary numbers. Nevertheless, the curve really does exist in this
space; it is just that the equation had to be put into a form that the graphing calculator
would recognize.
The Rose Curves: The final class of polar curve discussed in this section is the rose
curve. As suggested by its name, a rose curve consists of petals that emanate from a
single point. Rose curves has the general form r ( ) = a cos(b ) or r ( ) = a sin(b ),
where a and b are constants. The constant b is rather significant because it determines
how many petals the rose curve will have. If b is an even number, the number of petals
is 2b, while if b is an odd number, the number of petals is simply b. For instance,
consider the rose curve represented by r = 2 cos(6 ) :
Since b is an even number (6) the number of petals is twice b (12). What if b is an
irrational number like ? The number of petals will be irrational. However, if one
allows the limits on the magnitude of to be great enough, one would see a disc
generated. Consider the rose curve represented by r ( ) = cos( ) :
92
This is not a complete curve because the number of petals is irrational. The arrow
indicates where the curve stops. If one goes to WINDOW and makes max very large
(perhaps 20 ), the resulting graph looks a lot different:
If one thinks about this in terms of the philosophy of calculus, it is theoretically the
case that an irrational number as b results in an infinite number of petals. If one could
truly make = , the graph would look completely solid; it would be a complete disc.
Determining Points of Intersection: Solving polar equations is very important for the
next section of this chapter. It is not terribly difficult, but it requires more thought than in
solving equations in Cartesian or parametric coordinates. In general, one can decide
where two polar curves intersect by setting their polar equations equal to each other and
solving for . However, recall the property of polar coordinates that a -value
associated with a point is not unique; there are infinitely many other -values that are
associated with this same point. Thus, just because the -value of one polar curve at a
certain value of r is not the same as the -value of another polar curve at the same value
of r does not mean that the curves do not intersect. Setting the equations of the two polar
curves equal to each other yields the points where the curves intersect due to having the
same r- and -value. However, consider the two curves r1 ( ) = sin and
r2 ( ) = 1 sin . From the diagram below, there seems to be three points of intersection:
93
Setting the two equations equal to each other, one can find the -values at which the
curves intersect:
1
r1 ( ) = r2 ( ) sin = 1 sin 2 sin = 1 sin =
2
1 5
= , .
6 6
Thus, the two coordinates (r , ) where the two curves intersect are:
1
1 5
, and , .
2 6
2 6
But did it not seem, by looking at the graph, that there were three points of
intersection? If one were to zoom in on the graph on the TI-84, it would be clear that
both curves occur at the origin and, thus, intersect. Thus, there are three points of
1 1 5
intersection: , , , , and (0,0). Unfortunately, solving a trigonometric equation
2 6 2 6
will not confirm this because the solutions to these equations only show the points at
which both components of the coordinates are equal, and, as it turns out, while the rvalue of the coordinate of each curve at the origin is the same (0), the -value differs.
Solving for when r = 0 in each equation,
0 = sin = 0
0 = 1 sin =
Thus, at the origin, the coordinates for r1 are (0, 0) , while for r2 they are 0, .
2
However, these both represent the same point, the origin! The main message to take
away from this is that solving the resulting trigonometric equation from setting two polar
curves equal to each other might not convey all of the points of intersection because there
could very well be points of intersection in which the coordinates of the two curves are
not the same. The best way to check if this is the case is to analyze the actual graph on
the calculator.
The Calculus of Polar Curves
In this section, one will see a number of similarities between polar calculus and
parametric calculus. While the methods of differentiation and integration may seem
different from those encountered before, they still have the same interpretations and
applications; the derivative still refers to the slope of the tangent line and the integral is
still used to find area, arc length, or properties of solids of revolution.
Differentiation of Polar Equations: The formula for the derivative of a polar equation is
derived in the same manner as that applied to parametric equations; the chain rule is used.
94
dy
The derivative still has Leibnizs notation i.e. , but equations in polar form have x
dx
and y (together as r) in terms of . Thus, one applies the chain rule for derivatives:
dy
d
(r sin )
dy dy d d
d
=
=
=
.
dx
d
dx d dx
(r cos )
d
d
Applying the product rule to the numerator and denominator of the final expression
above:
d
dr
(r sin ) r cos +
sin
dr
d
d
is easy to evaluate
. Note that the expression
=
dr
d
d
(r cos )
cos r sin
d
d
because it represents the derivative of the polar curve with respect to , which does not
dr
does not represent the
require any algebraic manipulation. Be wary, however, that
d
dy
represents this. The AP Calculus BC
slope of the tangent line to the polar curve; only
dx
exam often asks questions pertaining to vertical and horizontal tangents in regards to the
dy
dy d
differentiation of polar curves. With the formula
in mind, a polar curve will
=
dx dx
d
dx
dy
dy
= 0 and
0 and a horizontal tangent when
=0
have a vertical tangent when
d
d
d
dx
and
0. Note that if both derivatives are zero in any of these cases, an indeterminate
d
0
form of results and no conclusion can be drawn.
0
Ex.) Find the horizontal and vertical tangents of to the limaon represented by
r ( ) = 3 4 cos on the interval [0, 2 ].
To find any vertical tangents, one sets
dx
dy
to zero, also making certain that
does
d
d
= 0,1.863996, , 5.0967858, 2 .
One must now make certain that these values do not yield a value of zero for
dy
:
d
dy
d
dr
r sin = r cos +
=
sin = 4 sin 2 3 cos 2 3 cos . Testing all of the
d d
d
dy
values of found above, one would find that none of them make
zero. Thus, one
d
can be certain now that there are vertical tangents at = 0,1.863996, , 5.0967858, 2 .
To find any horizontal tangents, one would set
dy
equal to zero and make certain that
d
dx
does not also equal zero:
d
dy
= 4 sin 2 3 cos 2 3 cos = 0. Again, one should solve this equation via
d
calculator by graphing in FUNCTION mode and finding where the graph crosses the xaxis on the interval. Doing this, one should get: = 1.6771037, 4.6060816.
dx
. Testing the two
One must now make certain that these values do not also make
d
dx
dx
, it is found that none of them make
zero. Thus, one can now be certain
values in
d
d
that there are horizontal tangents at = 1.6771037, 4.6060816.
One could check these answers by analyzing each of the discovered -values on the
actual graph of the polar equation.
Integration of Polar Equations: This section will cover the same three applications of
polar equations as was discussed in the section on parametric equations area, arc length,
and surface area of a solid of revolution. However, only the first two are tested on the AP
Calculus BC exam in regards to polar equations.
Area Under and Between Polar Curves: Due to the special characteristics of polar
coordinates in comparison to Cartesian coordinates, the formulae for area under and
between polar curves must be derived geometrically. Consider a small section of a polar
curve as shown below:
96
The diagram above represents a polar curve r = f ( ) that has an area bounded by the
angles and . Within this area is a circular sector with length r (or f ( ) ) , which
sweeps out the angle . Recall from circle geometry that the area of a circular sector
1
has the following formula: A = r 2 , where r is the radius and is the angle swept out
2
in radians. In accordance with the central theme of calculus (see the prelude), one moves
to the level of the infinitesimal and allows the number of sectors within this area to
become infinitely small, such that when these very small sector are added together, the
result is the value of the bounded area:
i =n
A = lim
i =1
1 2
1
r i =
2
2
r 2 d .
This integral does not look very complicated, and, indeed, it is not. However, perhaps
the most complicated part of solving problems concerning the area under a polar curve is
the determination of the limits of integration. Unlike the very simple determination of
the limits of integration in Cartesian coordinates, one must often closely analyze the
graph or solve trigonometric equations when dealing with polar curves. The following
several examples should clarify these necessities.
Ex.) Determine the area enclosed in the cardioid represented by r ( ) = 5(1 + sin ) .
The first step in any problem dealing with area under a polar curve should be an
analysis of the graph:
At the origin, one starts at = 0 . The origin is reached again once the directed
distance r has gone through one full revolution, or 2 . Thus, the limits of integration are
0 and 2 . One can now use the formula for area:
1 2
1 2
2
A=
(5 + 5 sin ) d =
sin 2 + 10 sin + 25 d . One can use the
2 0
2 0
1 cos 2
:
2
1
1 cos 2
+ 10 sin + 25 d =
2
2
1
51 1
51
97
51
80.11061267. Remember, if using the calculator, always find the area by going
2
to FUNCTION mode.
=
Ex.) Determine the area enclosed by the inner loop of the limaon represented by
r ( ) = e cos .
This problem is slightly more difficult. It requires a close analysis of the way in
which the curve is graphed. One must ask, At which -value does the loop begin? and
At which -value does the loop end? These questions can be answered by observing
the polar curve while it is being graphed. To ease this observation, one can slow down
the rate at which the calculator graphs the curve by decreasing the step size in the
WINDOW option. Zooming in on the part of the curve containing the loop,
The line indicates where the loop begins and the arrow indicates the direction in
which the curve is graphed. The loop starts at = 0 and that same point is reached when
= 2 . Unfortunately, 2 cannot be used as the upper limit of integration because that
would imply that one is taking the area of the entire limaon and not just the loop! To
find the area, one must exploit the property of symmetry in the above graph. Instead of
finding the area of the entire loop, one can find the area of one-half of the loop and then
multiply this value by 2. This is the only way to account for only the loop and not the
rest of the graph. Watch the calculator graph the curve again. The calculator finishes
graphing the first half (the bottom half) of the loop when it reaches the origin. To find
the -value associated with the origin, one plugs in 0 for r in the original equation and
solves for the appropriate value of :
e
(0) = e cos = cos 1 = 0.5251355796.
Thus, the bottom half of the limaon is completed when = 0.5251355796. One can
now find the area of the while loop:
Atotal = 2 A1 / 2loop =
0.5251355796
(e cos )2 d = 0.04985493.
Ex.) Determine the area of one petal of the rose curve represented by r ( ) = 3 sin 5 .
For this problem one must, again, consider the pattern of graphing and the symmetry
exhibited by the curve. The rose curve is shown below:
98
The calculator, of course, begins at = 0. The arrow indicates the direction is which
the calculator graphs the curve. Thus, the calculator finishes graphing the first petal
when it reaches the origin the second time. The -value at which this occurs can be
found by plugging 0 in for r in the original equation and solving for the appropriate value
of :
(0) = 3 sin 5 . This equation will equal zero when 5 = . Thus, =
. Now that
5
the two limits of integration have been found, one can find the area enclosed in one petal
of the rose curve:
1
A=
2
1
(3 sin 5 ) 2 d =
2
9 1 cos10 5 9 9 81
9 sin 2 5 d =
=
12.72345025.
=
22
2 0 2 10
20
It is often necessary not only to find the area under one polar curve, but the area
enclosed by two polar curves. In order to solve the integral that represents this area, one
must find at which -values the two curves intersect, for these will be the limits of
integration. The integral is simply the simply the difference between the definite integral
of the outer curve and the definite integral of the inner curve. For instance, in the
diagram below,
Ex.) Determine the area outside the circle represented by r ( ) = 2 and inside the
lemniscate represented by r 2 = 3 cos .
One must observe the graph to see which area must be found:
99
The shaded area represents outside the circle and inside the lemniscate. To
determine the limits of integration, one must find the -values at which the two curves
intersect:
2
r2 = r1 3 cos = 2 3 cos = 2 cos = = 0.8410686706. Note that
3
at this value of , only 1/4 of the shaded area has been completed. Thus, the difference
between the definite integrals of the two curves must be multiplied by 4 at the end:
A1 / 4 =
1
2
0.8410686706
dy 2
1 + dx . However, polar equations are
dx
dy
in terms
dx
of r and . The reader will recall from the subsection dealing with the differentiation of
dr
sin
r cos +
dy
d
polar equations that the result of this derivation was
. Substituting
=
dx dr
cos r sin
d
this expression into the formula for arc length,
expressed in terms of r and . Thus, one must use the chain rule to express
l=
2
dr
+
cos
sin
r
d
dx . For the sake of space, the algebra leading to a
1 +
dr cos r sin
d
100
l=
dr
r +
d .
d
2
Ex.) Determine the length of the first two turns of the logarithmic spiral represented by
3
r ( ) = e .
While the limits of integration could be found via calculator, trigonometric intuition
should convey that two turns means two revolutions, which means 4 . Thus, the
formula for the arc length of this spiral is l =
9
+ e 2 d = 20651.08.
16
Surface Area of Revolution of Polar Curves: One can derive the formulae for the surface
area of revolution of polar curves directly from the formulae for parametric curves:
dx dy
2 y + dt
dt dt
a
dx
dy
2 x + dt .
dt
dt
a
Recall that the radical expression is the arc length and that x = r cos and y = r sin .
Thus, the equivalent formulae for polar curves are:
dr
2 r sin r 2 +
d
d
dr
(When revolved about the x-axis): Surface Area =
2 cos r +
d .
d
Again, note that surface areas of revolution are not covered on the AP Calculus BC exam.
Ex.) Determine the surface area of the resulting solid when the entire rose curve
represented by r ( ) = 6 cos 2 is revolved about the x-axis.
This problem, like many problems involving integrals of polar equations, requires
the consideration of symmetry. This particular rose curve has 4 petals. One can, thus,
determine the surface area of revolution by revolving 1 petal about the x-axis and then
multiplying this result by 4. Notice where the calculator begins graphing the curve:
101
It begins at where the dot is indicated and then moves to the left, but does not return
to its original position until it has graphed the whole curve. Thus, an additional
complexity is added; one must determine the surface area of revolution of half of one
petal, multiply this result by 2, and finally multiply this result by 4. One can find the
upper limit of integration by determining the first -value at the origin:
(0) = 6 cos 2 . This expression is zero when 2 =
Surface Area1/2 Petal =
. Therefore, =
initial ray
pole
polar curves
Archimedean spiral
logarithmic spiral
limaon
cardioid
lemniscate
rose curve
102
103
Unfortunately, these vectors are not drawn exactly to scale; the vector for the
airplanes velocity would be longer than it is in the diagram if 30 km/h is actually
represented by the length it is given. Nevertheless, this diagram should convey the
essential features of vectors, that they depend both upon magnitude and direction.
Whenever describing a vector, one must always take both of these properties into
account. It is also important to note an important property of vectors; they can be moved
anywhere in space as long as their magnitude and direction are preserved. The positions
of the vectors used in the example are not unique. The smaller vector could have been
below the larger vector, the tail of the larger vector could have been drawn at the
head (i.e. the arrowhead) of the smaller vector, or any possible combination, as long
as both magnitude and direction are preserved.
Note that the notation for the vectors given in the previous example was quite
cumbersome. Scientists and mathematicians have developed more succinct ways to alert
a reader to a vector quantity, and sometimes, if one is only working in one dimension
(say, with a positive and negative axis), it usually suffices to describe the direction of the
vector quantity by its sign (i.e. negative or positive). For instance, suppose a particle is
only moving along the x-axis. This particles velocity, for instance, can be indicated by
the size of the number and a positive or negative sign, depending upon where the particle
is moving relative to the origin. In printed texts, vectors are usually indicated by boldface type. In handwritten problems, however, this is impractical, so one usually indicates
v
a vector by placing a small arrow above the symbol that represents a vector (e.g. V in a
handwritten problem would usually be seen as V in a printed text). More vector notation
will be discussed later in the chapter.
Vector Addition: Common operations upon vector quantities such as addition,
subtraction, multiplication, or division are carried out very differently from scalar
quantities. The primary reason for this is, of course, that one must take into account the
vector quantitys direction as well as its size. Addition and subtraction will be discussed
in this subsection on a case-by-case basis. In the context of vectors, subtraction is to be
considered a case of addition, in which one is adding a vector whose direction is
somehow denoted as negative. The term for the sum of vector quantities is the
resultant.
In one dimension, the addition and subtraction of vector quantities is relatively
straightforward. If, for example, two vectors are parallel to each other, meaning they
have the same direction in one dimension, then the vectors are simply added together and
the resulting vector has the same overall direction.
Ex.) Consider the two vectors below
5u
3u
where u represents some unit (e.g. m/s). Since these vectors have the exact same
direction in space, they are simply added together, yielding a resultant with the same
direction as the original vectors:
104
8u
In a slightly more complicated case, the vectors are antiparallel, meaning they are
oriented in exactly opposite directions in one dimension. In this case, one must specify
which direction, left or right, is positive and which is negative. In accordance with the
conventional Cartesian coordinate plane, the right is usually the positive direction and the
left is usually the negative direction.
Ex.) Consider vectors of the same magnitude as before, only now oriented in opposite
directions:
5u
3u
If the left direction is considered the negative direction, then one is essentially adding
3 and -5. The resultant vector is:
2u
Notice how the vector of larger magnitude was more important in the sense that it
determined the direction of the resultant vector.
In more complicated cases, vectors lie in mores than one dimension. In the case of two
dimensions, one must apply laws of geometry and trigonometry. Furthermore, while
angles of either 0 (for parallel vectors) or 180 (for antiparallel vectors) were implied in
the examples concerning vectors in one dimension, angles must be explicitly indicated in
two dimensions. The following examples discuss common situations of vectors in two
dimensions.
Ex.) Suppose one were asked to determine the resultant vector of the following two
vectors:
4u
6u
How would one begin to solve this? First, one must exploit a fundamental property of
vectors that has already been discussed; vectors can be moved anywhere in space as long
as their magnitude and direction are preserved. Thus, to facilitate this problem, one
positions the vectors like so:
6u
4u
105
The resultant vector connects the tail of the smaller vector to the head of the larger
vector to form a right triangle:
6u
4u
The magnitude of the resultant can be found by using the Pythagorean theorem and the
angle at which this vector is displaced upwards from the horizontal can be found through
right-triangle trigonometry:
R 2 = (4u ) 2 + (6u ) 2 R = (4u ) 2 + (6u ) 2 = 2 13 u
6u
= tan 1 0.9827937232 56.30993247 o.
4u
What if the vectors cannot easily be situated 90 apart to form two sides of a right
triangle? One method often used exploits the key properties of a parallelogram; opposite
sides and opposite angles are equal. In this so-called parallelogram method, the two
vectors in question are situated next to each other in a tail-to-tail fashion such that they
form the two sides of a parallelogram. To find the resultant vector, which is the diagonal
of the parallelogram, one would use the law of sines and cosines.
Ex.) Consider two vectors, one that lies on a horizontal axis and one that is displaced
40 from the horizontal axis:
5u
40
3u
Since vectors can be moved anywhere in space as long as their magnitude and direction
are preserved, one can form a parallelogram from these two vectors by situating
themtail-to-tail:
40
Given that, in a parallelogram, opposite sides and opposite angles are equal, the complete
parallelogram has the properties depicted in the following diagram:
106
One can now use the law of cosines to determine the magnitude of the resultant vector
(i.e. the diagonal of the parallelogram):
20
3u
20
5u
140
This vector is actually equivalent to the sum of two vectors, one that lies on the xaxis and one that lies on the y-axis. The magnitudes of these vectors are found by using
the following formulae:
Ax = A cos , Ay = A sin , where Ax and Ay represent the x- and ycomponents of vector A, repectively, A represents the magnitude of vector A (i.e., just
107
the number, which could also be represented by an unbolded A), and is the angle from
which vector A is displaced from the x-axis. Notice how the components of the vector
are not vectors themselves! This is because the formula from which they are derived
comes from right-triangle trigonometry, in which only magnitude is important. These
formulae should look familiar as they are essentially a conversion from polar coordinates
to Cartesian coordinates, as was discussed in the previous chapter. Vector resolution
represents both a mathematical and physical reality. Mathematically, resolving a vector
into its x- and y- components is basically an anti-sum, since one is finding the vectors
that could be added to yield the vector that one is resolving. This process also has
physical significance. Suppose someone is pulling a wagon with a force of 20 newtons
(N) at an angle of 60 from the horizontal. This person is actually pulling the wagon
forward and pulling the wagon up at the same time! Specifically, the wagon is pulled
forward with a force of (20 N) cos 60 o = 10 N and pulled upward with a force of
(20 N) sin 60 o 17.32050808. This person is actually pulling the wagon up to a greater
extent than he or she is moving it forward!
Vector resolution, as said before, is also useful in determining the resultant of three or
more vectors in two-dimensional space. In this process, one resolves each vector into its
x- and y-components (using the convention of measuring angles from the positive x-axis),
adds all of the x-components and all of the y-components, determines the magnitude of
the resultant vector ( R ) by converting from Cartesian coordinates to polar coordinates
2
with the formula R = R x + R y , and finally determines the angle at which this
resultant is displaced from the positive x-axis through the use of the formula
Ry
.
= tan 1
R
x
Ex.) Determine the resultant vector of the four vectors depicted below:
Start by resolving each vector into its x- and y-components, measuring all angles
from the positive x-axis:
A x = A cos = (10u ) cos 72 o = 3.090169944u
A y = A sin = (10u ) sin 72 o = 9.510565163u
108
One can now find the x- and y-components of the resultant vector, leading to the
magnitude of this vector:
R x = A x + B x + C x + D x = (3.09016994u ) + (12u ) + (12.28728066u ) + (7.250462296u )
= -4.447573016u
R y = A y + B y + C y + D y = (9.510565163u ) + (0u ) + (8.603646545u ) + (3.380946094u )
= 4.287864712u
R =
the negative sign associated with the x-component R was not included because this
formula is derived from the Pythagorean theorem, and negative signs have no meanings
in classical geometry) The direction of R can be determined in the next step. Note that
the x-component of R will not be shown as negative here as well, since the following
formula is derived from right-triangle trigonometry, and triangles cannot have negative
dimensions:
Ry
4.287864712u
o
= tan 1
= tan 1
= 43.95258913 .
4.447573016u
Rx
Where exactly is this angle (i.e. in what quadrant)? To remove this ambiguity, one
examines the signs of Rx and Ry. Since the former was originally found to be negative
and the latter is positive, this angle lies in quadrant II. Thus, when measured from the
positive x-axis, the angle is 180 o 43.95258913o = 136.0474109 o. As a final answer, the
resultant vector is 6.177919514u, 136.0474109 from the positive x-axis.
This was quite a bit of information to deal with, so here is an overview of the theory
behind it:
Any vector can be resolved into its components, which are not truly vectors
themselves, due to their geometric and trigonometric derivation. Vector resolution is
particularly helpful in determining the vector sum (i.e. the resultant) of three or more
vectors in two-dimensional space. In this process, one first determines the x- and ycomponents of all the vectors in question. Note that these components do not
intrinsically have direction, though they may be positive or negative depending upon
their angle from the positive x-axis. One then sums all of the x-components and all of the
y-components, keeping the negative signs in the process. One then converts from
109
Cartesian coordinates to polar coordinates, using only positive values for Rx and Ry. Once
both the magnitude and associated angle of the resultant are found, one must determine in
which quadrant it lies. One does this by examining the original values of Rx and Ry.
There is a good intuitive way to informally check ones answer to a problem of this
type. Imagine that the vectors in question all represent forces tugging on something at
the origin. The resultant vector determines in which direction that thing will move.
Based upon the lengths and directions of the arrows representing the vectors, one should
be able to tell which magnitude and direction will win this strange game of tug-of-war.
In the examples already presented, the combined effects of vectors C and D are
qualitatively greater than those of A and B. Thus, the resultant vector will most likely be
pointed toward the negative x-axis, which was found to be true. Furthermore, the
combined effects of vectors A and D tugging upwards seems to overshadow the action
of vector C. One should, therefore, suspect that the resultant vector is pointed upwards,
which was also found to be true.
Unit Vector Notation and Generalization to Three Dimensions: There is a more concise
way to express the magnitude and direction of a vector than through the use of polar
coordinates. Indeed, an answer such as 6.177919514u, 136.0474109 from the positive
x-axis is rather burdensome when doing problems. Instead, one often wishes to express
a vector in terms of its components. However, as was discussed earlier, the components
of a vector are not vectors themselves, so simply displaying them does not suffice to
characterize a vector quantity. In order to use vector component notation and actually
describe a vector one must understand the concept of the unit vector. A unit vector is a
dimensionless (i.e. without units) vector that has a magnitude of unity (1 thus the name
unit vector) and that points along an axis. In two dimensions, a unit vector represented
by i points along the x-axis and a unit vector represented by j points along the y-axis:
How do unit vectors allow one to write a vector in terms of its components? By
multiplying these unit vectors to the components associated with the same axis, one can
easily express a vector by the sum of its components multiplied by these unit vectors. To
clarify, consider vector A below:
110
111
Notice that the extra dimension added is known as the z-axis. The unit vector
associated with this axis is denoted as k (or k if handwritten). Thus, one can now extend
the formulae for vectors in two dimensions to vectors in three dimensions:
A = Ax i + Ay j + Az k
A = A=
Ax + Ay + Az .
In the field of linear algebra, these formulae are extended to even more dimensions,
but that will not be discussed here.
Multiplication with Vectors: Before moving on to more difficult material in this
subsection, it is important to note some algebraic properties of vectors when
multiplication is concerned. If A and B are two vectors in three dimensions and m and n
are two scalars (i.e. with only magnitude, no direction), then the following algebraic
statements hold true:
mA = mAx i + mAy j + mAz k
(m + n) A = mA + nA = (mAx i + mAy j + mAz k ) + (nAx i + nAy j + nAz k )
n(mA) = m(nA) = mn Ax i + mn Ay j + mn Az k
The key point to take from these expressions is that when one multiplies a vector by
a scalar, nothing fancy is required; it suffices just to multiply every component of the
vector by the scalar. However, the same simplicity does not hold for multiplying a vector
by another vector. In general, there are two ways in which one can multiply two vectors.
One of these results in another scalar and one results in a vector.
The Scalar (Dot) Product: In general terms, the scalar product of two vectors in twodimensional space is described as the product of the magnitude of one vector and the xcomponent of the other vector. The scalar product is indicated by a dot () , hence the
alternative name dot product. The equation for the scalar product of two vectors A and B
separated by an angle is given as A B = A B cos = AB cos . The geometric
interpretation for this is shown below:
112
Thus, the scalar product is the product of the magnitude of a vector and the projection
of another vector onto that vector. While the name of this type of product certainly gives
it away, note that the result is a scalar quantity, not a vector quantity.
It is important to note special algebraic properties of the scalar product:
A B = Ax B x + Ay B y + Az B z This expression relates to vectors in R3. It basically
states that the scalar product is the sum of the ordinary products of the components in
each dimension.
A B = B A Scalar products are commutative.
A (B + C) = A B + A C Scalar products are distributive
m( A B) = (mA) B = (mB) A = m( A B) , where m is a scalar. This expression
extends the distributive law to scalars.
i i = j j = k k = 1, i j = i k = j k = 0. Recall that unit vectors have magnitudes of
unity. The first set of scalar products is equal to unity because they involve
multiplication of two vectors with the exact same direction (and magnitude). The second
set of scalar products is equal to zero because they involve multiplication by vectors that
are perpendicular. To see why this is the case, refer back to the diagram of the righthanded coordinate frame.
In physics, the most common application of the scalar product is the
determination of a quantity known as work (W), which is the scalar product of the force
(F) exerted on an object in newtons (N) and the displacement (s) (the change in position)
of that object in meters (m). Thus, W = F s = Fs cos . This equation only holds for a
constant force. Note that when the force and displacement vectors have the same
direction, the angle between them is 0. Thus, when the force is in the same exact
direction as the displacement, work is simply the product of F and s, since cos (0) = 1.
When the force vector is perpendicular to the displacement vector, no work is done on the
object in question because cos (90) = 0. Thus, while it may seem like a person would be
doing a great deal of work when lifting a cinder block while he or she is walking, this
person is doing absolutely no work in the physical sense!
Ex.) A wagon is pulled with a constant force of 25.0 N for a distance of 30 m as
shown by the diagram below:
113
114
This diagram represents the vector product A B . Notice that the fingers of the
right hand start at A and curl toward B. Since A came first in the equation, it cannot be
the other way around (i.e. one cannot start at B). The thumb points upward, so this is the
direction of the vector product A B .
This diagram represents the vector product B A . In this case, the fingers of the
right hand start at B and curl toward A. Since the thumb points downward, this is also
the direction of the vector product B A . Notice that this direction is the direct opposite
of A B . This supports the algebraic statement A B = B A . Note, however, that
both vector products are perpendicular to the plane in which the two vectors lie. In the
diagram, this was indicated by drawing in three dimensions. One could also use the
symbols and . The former indicates that the vector product is coming out of the
plane of the paper, while the latter indicates that the vector product is going into the
plane of the paper. Thus, the symbol would be associated with the first diagram and
the symbol would be associated with the second diagram. Note that there is a way to
determine the exact unit-vector notation of the vector product, but it involves a linearalgebraic tool known as the determinant, which will not be covered here.
Vector products are important in describing phenomena in many areas of physics.
Two applications will be discussed here. The first concerns rotational dynamics, or the
study of the forces that cause rotational motion. Everyone is probably familiar with the
fact that a door is more difficult to close when pushing closer to its hinges. However, one
can close the door almost effortlessly by exerting a force at the very end, away from the
hinges. Why is this the case? The force that one is exerting to cause rotational motion is
called the torque (represented by the Greek letter tau ). Torque is a vector quantity
that is directly proportional to the magnitude of the force applied and the distance from
the axis of rotation at which this force is applied. Thus, for a certain torque to be
achieved, if one starts at a small distance from the axis of rotation, a greater force would
need to be applied. Similarly, if one is farther away from the axis of rotation, one would
need to apply less force to achieve the required torque. Torque also depends upon the
angle at which the force is applied to the object in question. Applying the force
115
perpendicular to the object maximizes the torque while applying the force parallel to the
object causes no torque. Thus, a force that is applied perpendicular to a door will
maximize its rotational motion, while a force that is applied to the side of the door toward
the hinges will (obviously) cause no rotational motion at all. All of these elements can be
summarized by a vector product: = = F r = Fr sin , where F is the force applied,
r is the position relative to the axis of rotation at which the force is applied, and is the
angle between the force and position vectors.
Ex.) The following diagram represents a 3.00-m long wooden rod attached to a hinge
that allows the rod to rotate freely on its axis with negligible friction.
Unfortunately, this is not in a form conducive to the vector product. In the vector
product, the vectors are connected head-to-head like so:
A little geometric analysis would show that, luckily, the angle between these two
vectors in this situation is also 30. With this information, one can find the magnitude of
the position vector:
116
(6.75 N m)
= 0.27 m .
F sin (50.0 N) sin 30 o
Note that while it was not explicitly a part of this problem, the direction of the torque
would be into the page ( ) because the fingers of the right hand are starting at the
force vector (because it appears first in the equation) and curling toward the position
vector.
Another interesting application of the vector product is in the field of
electromagnetism. Electricity and magnetism, due to the work of the 19th-century
Scottish physicist James Clerk Maxwell, are united based upon the fact that electric
charges in motion generate magnetic fields. This is fundamentally true, even in, say, a
refrigerator magnet! It is also the case that a preexisting magnetic field will exert a force
upon a moving point of charge. This force is directly proportional to the magnitude of
the charge, the velocity of the charge, and the strength of the magnetic field. The force
exerted on the point of charge by the external magnetic field also depends on the angle at
which the charged particle moves relative to the magnetic field. The magnetic force is
maximized when the charge moves perpendicular to the magnetic field and is zero when
it is parallel to the magnetic field. All of these observations should signal that the force
exerted by a magnetic field on a charged particle can be represented by a vector product:
F = F = qv B = qvB sin , where F is the force of the magnetic field, q is the
= Fr sin r =
magnitude of the charge (Note that charge is always a scalar quantity), v is the velocity of
the charged particle, B is the strength of the magnetic field, and is the angle between
the velocity and magnetic fields vectors. Charge magnitude is measured in coulombs (C)
and magnetic field strength is measured in newtons per coulomb-meters per second
(N/Cm/s), which is succinctly referred to as the tesla (T).
Ex.) Consider a charged particle that is moving out of the page in a magnetic field that
is directed to the left:
This particle has a charge of 9.1314 10 19 C and is moving out of the page with a
velocity of 3.3017 10 2 m/s . If the strength of the magnetic field is 2.7877 T and the
charged particle experiences a force of 5.1860 10 16 N, what is the angle between the
velocity vector and the magnetic field vector?
F
(6.1860 10 16 N)
F
= sin 1
= sin 1
19
2
qvB
qvB
(9.1314 10 C)(3.3017 10 m/s)(2.7877 T)
= 47.393o.
117
Vector-Valued Functions
Note: This section is on the AP Calculus BC exam!
Thus far, the theory and application of vectors have been extensively discussed. This
section introduces perhaps an even more highly powerful use of vectors in mathematical
analyses. Scientists and mathematicians often express functions in terms of vectors to
analyze curves in space, and this will be the theme for the remainder of the chapter. A
vector-valued function, or simply vector function, is, not surprisingly, a function defined
by vectors. More formally, if real scalar values denoted by u can be located with a vector
A, then A is function of the scalar u: A(u). Suppose that the scalar values have x-, y -,
and z-coordinates. If this is the case, then A(u) can be expressed in terms of the
component vectors of A: A(u ) = Ax (u )i + Ay (u ) j + Az (u )k. This is the point where the
relationship between vector functions and parametric equations comes into play; each
component vector of the vector function is, itself, a parametric equation with u as the
parameter! To see why this is the case, consider a more familiar example from physics
the position function. This is actually a vector function denoted as r. A particles
position in three-dimensional space (i.e. x-, y -, and z-coordinates) with respect to the
time parameter t can be found through the vector function r (t ) = x(t )i + y (t ) j + z (t )k .
This can also be understood pictorially:
This diagram depicts the position vector r locating a particle in R 3 at a certain point
(x, y, z). It is important to note that this particles path is actually a parametrized curve.
The parametric equations are x(t ), y (t ), and z (t ). Thus, a vector function is essentially a
sum of parametric equations with the addition of unit vector notation. As stated in the
introduction to this chapter, this is basically the extent to which the AP exam tests the
theory of vectors.
The Calculus of Vector-Valued Functions
Note: This section is also on the AP Calculus BC exam! Note also: Now that calculus is
being applied again, be sure the graphing calculator is in radian mode!
The theorems of calculus as applied to vector functions are exactly analogous to those
of functions not defined by vectors. The only fundamental difference is that the calculus
theorems of vector functions include vector notation. For example, the limit definition of
118
the derivative is not necessarily unique for vector functions. For a vector function A(u ),
the limit definition of the derivative is as follows:
dA
A(u + u ) A(u )
= lim
. This is exactly analogous to the definition
du u 0
u
It is also important to know the following formulae in relation to particle motion in vector
space:
2
speed = v (t ) = v x + v y + v z
a x (t ) =
d v(t )
119
Integration of vector functions is also quite straightforward, for the rules of integration
are the same. Integrating a vector function means integrating all of the components of the
vector function. For a vector function A(u ),
C = Cxi + C y j + Czk .
t2
v (t )dt =
t1
t2
v x (t )dt +
t1
t2
t1
v y (t )dt +
t2
t1
v y (t )dt = r
Is there significance for the definite integral of the speed of a particle from one time to
another? To see if there is, one can examine the integral and see if it looks like anything
familiar. Suppose one is working in two dimensions rather than three. The definite
integral of speed would have the following appearance:
t2
t1
v (t ) dt =
t2
t1
v x + v y dt =
t2
t1
dy
dx
+ dt.
dt
dt
This last integral should look familiar; it is the formula for the arc length of a
parametric curve! Since this is the length of the particles path, the definite integral of
the speed of a particle from one time to another equals the total distance distance traveled
by the particle during that period of time. What is the difference between change in
position (displacement) and total distance traveled? Consider a particle moving along the
following path in two dimensions:
120
v (t )dt =
The integral of the x-component of velocity requires integration by parts, but since
one term is a polynomial, one can use the tabular method (see chapter 1):
u
dv
1
t
5t
+1
e
-1
5
e t
t
+1
0
e
-1
Thus,
(5t e
)idt =( 5t e
121
To find the y-component of the position, one must integrate the y-component of the
velocity:
3
3t 3 j dt = t 4 j + C y = y (t ) j + C y . Since the y-ccordinate of the particle at
4
3
t = 0 is 0, Cy is also 0, so y (t ) j = t 4 j . Thus, the position equation is:
4
3
r (t ) = 5t e t 5e t i + t 4 j.
4
b.) Recall that speed is the magnitude of the velocity vector. So, at t = 3:
speed t =3 = v (3) =
[5(3)e ] + [3(3) ]
( 3) 2
3 2
= 6561.558.
c.) The total distance traveled is the definite integral of speed, or the arc length of the
parametric curve, so:
total distance t =0,t =4 =
v x 2 + v y 2 dt =
t 2
dy
3(3) 3
= dt =
= 108.4618994.
dx t =3 dx
5(3)e ( 3)
dt t =3
x(3) = 0.9957414
y (3) = 60.75
y (60.75) = (108.4618994)( x + 0.9957414) .
The Differential Geometry of Curves in Space
122
The TNB Frame: In the study of curves in space, a generic vector function is usually
represented as the position function: r (t ) = x(t )i + y (t ) j + z (t )k. This function is
assumed to be differentiable. The material introduced in this subsection is most
popularly applied to motion of particles in space, especially as applied to flight.
Nevertheless, the position function is still useful in analyzing curves that do not model
motion. For instance, the double-helix model of a DNA molecule can be analyzed from
the standpoint of the position function as well.
In R 3 , curves are often analyzed through the use of three types of vectors. The
first is the unit tangent vector (T). This vector is useful because it quantifies the
direction in which a particle is moving at a certain point on a curve. A special property
of the unit tangent vector is that it is magnitude-blind; it provides a measurement solely
v
for direction. This is evident in the formula for the unit tangent vector: T = . Notice
v
how division by the magnitude of the velocity vector (speed) allows the unit tangent
vector to solely quantify direction. Note also that this division results in a magnitude of
unity for this vector. The unit tangent vector is also sometimes conveyed as the
derivative of the position function with respect to arc length (which, in this case, is
dr
represented by s). Thus, it is also the case that T = . This formula is very impractical
ds
because it is quite difficult to find position in terms of arc length. The second vector used
to analyze curves in space is called the principal unit normal vector (N). This vector is
orthogonal (i.e. perpendicular) to the unit tangent vector. The principle unit normal
vector primarily conveys the extent to which a particles path deviates from a straight
line. Its formula is derived from a scalar quantity known as curvature ( the Greek
letter kappa). Curvature (besides the obvious definition of extent to which a curve
curves!) is related to the rate at which the curve changes direction at a point. Thus, it is
defined mathematically as the magnitude of the derivative of the unit tangent vector with
dT
respect to arc length: =
. Again, it is impractical to manipulate a function to be
ds
defined by arc length. Thus, one can use the chain rule in the following way:
dT ds
=
, where the latter term is simply the reciprocal of speed. It is also possible to
dt dt
simplify the expression of even more by expressing it in terms of familiar aspects of
motion. While the derivation will not be shown here, curvature can also be found from
va
the following formula: =
(Note that the numerator is a vector product). How does
3
v
curvature relate to the principle unit normal vector? By multiplying the derivative of the
unit tangent vector with respect to arc length by the reciprocal of curvature, one yields the
dT 1
principle unit normal vector: N =
. Unfortunately, this equation is in terms of arc
ds
length, so the chain rule is used to yield a formula in terms of t:
123
dT dt
dT
1
d T 1 dT
=
= dt ds = dt . Due to this division, the principal unit
N=
dT dt
dT
ds ds dT / ds
dt ds
dt
normal vector also has a magnitude of unity. The third vector used in curve analysis is
the unit binormal vector (B). This vector is orthogonal to both the unit tangent vector
and the principal unit normal vector. Loosely speaking, the unit binormal vector
measures the extent to which a moving particle is dissociating itself from its path. The
unit binormal vector is simply the vector product of the previous two vectors discussed:
B = TN .
It is now time to put all of this information together. The unit tangent vector, principle
unit normal vector, and unit binormal vector all have unit length and lie orthogonal to
each other. Thus, they form a three-dimensional coordinate frame, known as the TNB
frame, or Frenet frame, in honor of the 19th-century French mathematician Jean Frdric
Frenet (pronounced fruh-nay):
This frame is always associated with the particle in question. That is, as the particle
moves along a curve in space, the TNB frame always moves along with it. Throughout
the particles journey, the TNB frame describes different aspects of that particles
motion; T describes the direction of the particle at a particular point, N describes the
extent to which that particles path deviates from a straight line, and B describes the
extent to which a particle could potentially take another path besides the one it is
actually taking. It may be more helpful to understand these aspects of particle motion
from a common experience (for some) riding on a roller coaster. In this case, T would
represent the direction tangent to the coaster one feels that he or she is moving, N would
represent the extent to which one feels like he or she is curving, and B would represent
the extent to which one feels as if he or she is being pushed up. The following diagram
represents a particle traveling along a curve in space and the TNB frame associated with
it at specific points on the curve:
124
There is one more useful aspect of particle motion on curves in space to be discussed
torsion ( ). Torsion is a scalar quantity that describes the extent to which a curve
twists in space. Torsion is equal to the magnitude of the rate at which the direction of
dB
the unit binormal vector changes with respect to arc length: =
. There is a way to
ds
simplify this formula, but it involves the determinant, which is beyond the scope of this
book..
As one can glean from the formulae in this subsection, problems involving the TNB
frame are lethally tedious, and it is almost always the case that even the most
computationally careful of humans will make some sort of error. Thus, in the upcoming
example, the calculations will be done on Mathematica.
Ex.) A malfunctioning model airplane moves through space as described by the position
equation: r (t ) = sin(3.33t )i + cos(2.22t ) j + 1.11(t )k.
a.) Determine the curvature of the airplanes path at three times:
t = 0, t = , and t = 2 .
b.) Compare the equation of the derivative of the unit tangent vector with respect to
time to the equation for curvature. What is being compared?
It would be helpful to visualize this curves path using Mathematica :
ParametricPlot3D[{Sin[3.33*t],Cos[2.22t],1.11*t},{t,0,2*Pi}]
-0.5
5
-1
0.5
0
-1
1 -0.5 0 0.5
Graphics3D
Indeed, this plane is experiencing some problems. a.) To compute the curvature at the
three desired times, one uses the [r_][t_] function in Mathematica :
[gamma][0] = 0.4,[gamma][Pi]= 1.52186,[gamma][2*Pi]=0.897866
125
Note that [gamma]merely stands for the position function, for the purpose of
saving space. This plane is not even screwing up consistently! Its curvature increases
and then decreases on the time interval.
b.) For this part of the problem, one can use the same Mathematica operator,
but leave the result in terms of t. For instance, the curvature function can be found in the
following way:
,
2
2
I H29.9266 Cos@2.22 tD + 151.504 Sin@3.33 tD +
[gamma][t]=
2 32
[r (t )] =
D[{UnitTangent[gamma][t]},t]
11.0889 Sin@3.33 tD
!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!
!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!
1.2321 + 11.0889 Cos@3.33 tD2 + 4.9284 Sin@2.22 tD2
32
!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!
!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!
1.2321 + 11.0889 Cos@3.33 tD2 + 4.9284 Sin@2.22 tD2
>>
32
Notice how this expression, which is even more confusing, gives three
coordinates instead of a single equation. This is the case because this is the derivative of
the unit tangent vector with respect to time, while the curvature was the derivative of the
unit tangent vector with respect to arc length.
126
assigned to each spatial point is temperature in degrees Celsius. Assume that in the
region being studied, temperature only varies in the x-direction:
If T is a function of x, y, and z and temperature only varies in the x-direction from T1 and
the vector u indicates the direction from T1 to T2 then the directional derivative
suggested by the diagram is:
T T
T
u = 2 1 u.
x y , z
x 2 x1
While the directional derivative in this case was chosen to be from T1 to T2, it could have
been from T1 to any other temperature. Thus, given a point in space that represents a
scalar quantity, there are an infinite number of directional derivatives from that point.
One directional derivative is especially significant. It is called the gradient. The gradient
is a vector that points in the direction of greatest increase in the scalar quantity in
question. The gradient is often denoted by a symbol, which is referred to as the del
operator. The gradient of a scalar field function f is defined mathematically as:
f
f
f
f =
i+
j + k.
x
z
dy
Gradients appear extensively in the physical sciences. For instance, the one-dimensional
dU
equation that relates force ( Fx ) to potential energy (U) is Fx =
. This can be
dx
U U U
+
+
.
extended to three dimensions through the use of the gradient: F = U =
x
z
y
In electrostatics, the gradient relates the electric field (E), a vector, to electrical potential,
V V V
or voltage (V), a scalar: E = V =
+
+
.
x y z
Differentiation in Vector Fields: Differentiation in scalar fields requires at most three
partial derivatives, one for each space component. Differentiation in vector fields,
however, requires more partial derivatives because both elements of the derivative are
vectors, each with different components. Consider differentiation in a force field
F( x, y, z ) . Unlike differentiation in a scalar field of temperature, the force vector itself
has three components; Fx , Fy , and Fz ; each of which are a function of x, y, and z. Thus,
for differentiation in a three-dimensional vector field, there are a total of nine different
partial derivatives! It is quite impractical to describe a physical system in this manner.
Thus, two new mathematical expressions are introduced that more concisely describe
128
spatial variation in a vector field. The first is divergence, which, for a vector field
function F, is often denoted as divF . Physically speaking, divergence measures the
extent to which a point in space is a source or sink for something, making it an ideal
quantity in physical systems characterized by the flow of matter and energy.
Mathematically, divergence is the scalar product between the del operator and the vector
field in question. For a vector field F:
Fy Fz
F
divF = F = i +
j + j ( Fx i + Fy j + Fz k ) = x +
. Note that
+
y
z
x
y
z
x
divergence is a scalar quantity. The second mathematical expression often used to
describe spatial rates of change in vector fields is called curl, which, for a vector field
function F, is often denoted as curl F. Note that curl is bolded because it, unlike
divergence, is a vector quantity. In physical terms, curl can be described as the extent of
rotation of a vector field, making it an ideal mathematical expression to use in fluid
dynamics, such as studying vortices. While divergence is the scalar product between the
del operator and the vector field in question, curl is the vector product between these two
elements (The last part of this expression was arrived at through the determinant):
F Fy Fx Fz Fy Fx
i +
k.
curl F = F = i +
j + j ( Fx i + Fy j + Fz k ) = z
j +
y
z
y
z
z
x
x
y
The Line Integral: This last subsection introduces an important mathematical expression
in vector calculus the line integral. This kind of integral, also sometimes referred to as
a path integral, measures scalar or vector quantities along a curve in space. One of the
most popular examples of a line integral concerns work, a physical quantity that was
introduced in the first discussion of the scalar product in this chapter. This section,
however, will consider the work done by a non-constant force. Consider a particle
moving along a path C described by the position function r. This particle is acted upon
by a non-constant force F(r) along this path. Recall that the work done on an object is
the scalar product of force and displacement. This formula still holds here, but since a
non-constant force is doing the work, one must move to the differential level and divide
the particles path into an infinitely large number of infinitesimal pieces in which the
force is constant. Taking the sum of all of these instances of work yields the actual
work done by the non-constant force. This, of course, involves an integral, specifically a
line integral: W =
F(r) dr. Note the new notation of this integral. It is not a definite
C
integral, but the subscript C signals that this integral is being evaluated over a curve. In
order to evaluate such an integral, it is necessary to manipulate it to yield an equivalent
definite integral. In realizing that the vector function r is actually a parametric equation
with the time parameter t, one can also express the force equation in terms of t:
F(r ) = Fx (t )i + Fy (t ) j + Fz (t )k. Taking the scalar product of the parametrized force and
differential position functions yields the following:
dy
dx
dz
129
W =
F(r ) dr =
t2
t1
dy
dz
dx
Fx (t ) + F y (t ) + Fz (t ) dt .
dt
dt
dt
While this line integral involved motion over a so-called open path, line integrals can
also be evaluated over closed paths, in which the initial point coincides with the endpoint.
For instance, consider the closed path below:
If a non-constant force is acting upon a particle traveling along this curve, the notation
of the line integral must indicate that the path is closed. This is done by the inclusion of a
small circle around the integral symbol:
on the particle can be expressed as the sum of the line integrals from point a to b to c and
point c to d to a:
W =
F(r ) dr =
abc
Oftentimes, a line integral only depends upon the points along the curve being
considered, not on the actual path taken. For instance, consider three possible paths from
point A to point B below:
While each path is rather distinct from the others, it may very well be the case that the
line integral is the same for each path. In this case of path independence, only the initial
and final points have any importance in the value of the line integral, not the actual path
taken. If this is the case, the field in question is called a conservative field. For instance,
the work done by gravity is irrespective of the path of an object,t while the work done by
friction does depend upon the path. Thus, gravity is an example of a conservative force
and friction is an example of a non-conservative force.
130
Concluding Remarks
This final chapter introduced the theory and calculus of vectors, a very important
concept in mathematics and science. The nature of vectors and the algebra that reflects
this nature were discussed. After vectors were considered from this relatively unique
standpoint, they were then discussed in the context of functions, which conveyed their
close relationship with the parametric coordinate system, especially in terms of position
velocity and acceleration. It was discussed that the theorems of vector calculus are
essentially the same as those for scalar functions, and this calculus was applied to
problems in motion. The final two sections introduced concepts that are generally
introduced in Calculus III. The penultimate section introduced the theory of the
differential geometry of curves and the mathematical expressions used to analyze curves
in space. The final section heavily converged upon multivariable calculus in discussing
the calculus of scalar and vector fields.
Key Terms:
vector quantities
scalar quantities
vector
resultant
vector resolution
unit vector
right-handed coordinate frame
scalar product
work
vector product
right hand rule
torque
vector-valued function
differential geometry
unit tangent vector
principal unit normal vector
curvature
unit binormal vector
TNB frame
torsion
field
scalar field
vector field
field functions
scalar field function
vector field function
directional derivative
gradient
del operator
divergence
curl
line integral
conservative field
131
132
Practice Exam
Section 1, Part A
55 minutes
28 questions
Directions: Solve each of the problems in this section, using the available space for
scratchwork. No credit will be awarded for anything written in the test booklet. Use
your time wisely. No calculator of any kind may be used in this section.
Note: Unless otherwise indicated, the domain of a function f is the set of all real
numbers x for which f (x) is a real number.
1.) If 3xy = x 4 y 3 , then
a.)
dy
=
dx
4x 3 y 3 3 y
3y 4 y3 x3
3x 4 y 2 + 4 x 3 y 3
, c.)
, b.) 4 2
, d.)
3
3x 3x 4 y 2
x 3y 3
9x 4
2x8
3x
x4
, e.)
9x 4
2 yx 8
________________________________________________________________________
2.) The growth of a population is modeled by the differential equation
dN N
N
= 1 , where N is the number of individuals in the population. How many
dt 12 72
individuals will there be when the carrying capacity is reached?
72
________________________________________________________________________
133
3.)
3 sin xdx =
1 1
1
cos t , d.) 3 cos t , e.) cos t
3 3
3
________________________________________________________________________
4.) Consider the area enclosed by the two curves shown below:
The formula for the volume of the solid generated by revolving the shaded area about
the y-axis is
5. 5
b.) 2 ( g ( x) f ( x) )xdx
c.)
[1 (g ( x) f ( x))] xdx
d.)
[1 + (g ( x) f ( x))] ydy
e.) 2
(g ( x) f ( x))ydy
a.)
(g ( x) f ( x) )2 dx
5. 5
5. 5
f (5.5)
f (0)
f ( 5. 5 )
f (0)
134
________________________________________________________________________
5.) A particle moves along a path described by the parametric equations
x(t ) = 5 sin t and y (t ) = t 3 . At t = 4, the particle leaves the path and begins to travel along
the tangent line to the path at that point. What is the slope of this line?
a.)
3 2
5
5 cos t
5
3
, e.) t 2 sec t
t sec t , b.) , c.) , d.)
2
5
3
3
5
3t
________________________________________________________________________
6.)
a.)
b.)
c.)
d.)
e.)
2x3 + 2x 2 + 1
dx =
x 4 + 2x3 + x
1
ln x 4 + 2 x 3 + x + C
6
1
ln x 4 + 2 x 3 + x + C
2
3
2 ln x + ln x 3 + 2 x 2 + 1 + C
4
1
ln x + 3x 3 + 2 x + C
2
1
1
2 ln x + x 3 + x + C
3
2
________________________________________________________________________
7.) Which of the following series converge?
I.)
n =1
n2 + 3
, II.)
n2
n =1
2n!
, III.)
10 n
n =1
3n
4n + 1
a.) I only, b.) II only, c.) III only, d.) Both I and III, e.) Both II and III
135
________________________________________________________________________
8.) If y = 3 ln(3 x) + 5 , then y =
a.)
b.)
c.)
1
3 3 [ln(3x) + 5] 2
1
3( x + 5) 3 [ln(3x) + 5] 2
1
3x 3 [ln(3 x) + 5] 2
d.) 3x 3 [ln(3x) + 5] 2
e.) 3 3 [ln(3x) + 5] 2
________________________________________________________________________
9.)
a.)
b.)
c.)
d.)
e.)
3 x 3 cos x dx =
3 4
x sin x + C
4
3x 3 sin x + 9 x 2 cos x + C
3x 3 sin x + 9 x 2 cos x 18 x sin x 18 cos x + C
3x 3 sin x 9 x 2 18 x sin c 18 cos x + C
3 4
3 5
3 6
3 7
x cos x +
x sin x
x cos x
x sin x + C
4
20
120
840
136
________________________________________________________________________
10.) The function f (x) is defined piecewise as:
2
f ( x) = x 1, x < 2
3a 2, x 2
1
3
1
5
3
, b.) , c.) , d.) , e.)
2
4
3
3
5
________________________________________________________________________
11.) The graph below represents f (x) , a functions derivative.
On the x-interval shown, what is the total number of minima and maxima on the graph
of the function f (x) ?
a.) 5, b.) 6, c.) 7, d.) 8, e.) Not enough information provided.
________________________________________________________________________
137
1
12.) lim 2 =
________________________________________________________________________
13.) Find the value of x that satisfies the Mean Value Theorem for f ( x) = x 3 + 2 x + 1 on
the closed interval [0, 2].
a.) 1, b.) 6, c.)
2
3
, d.)
3
, e.)
2
________________________________________________________________________
138
14.)
a.)
dx
4 + 9x 2
1
3
3
1
________________________________________________________________________
n =0
series for xe
( x +1)
xn
, which of the following is a power
n!
( x + 1) 2 ( x + 1) 3
+
+L
a.) 1 + ( x + 1) +
2!
3!
x( x + 1) 2 x( x + 1) 3
+
+L
b.) x + x( x + 1) +
2!
3!
( x + 1) 3 ( x + 1) 4
c.) ( x + 1) + ( x + 1) 2 +
+
+L
2!
3!
x4 x5
d.) 1 + x 3 +
+
+L
2! 3!
x( x + 1) 3 x( x + 1) 4
2
e.) x + x( x + 1) +
+
+L
2!
3!
139
________________________________________________________________________
Questions 16 and 17 refer to the graph below:
16.) The above diagram is a graph of the function f (x) on the interval [-2, 10]. If
g ( x) =
________________________________________________________________________
17.) Where on the interval [-2, 10] is g (x) increasing?
I.) [-2, 0], II.) [0, 3], III.) [3, 5], IV.) [5, 10]
a.) I only, b.) II only, c.) I and III, d.) II and IV, e.) IV only
________________________________________________________________________
18.) A spherical balloon is inflated with helium at a rate of 200 ft3/min. How fast is the
balloons radius increasing at the instant the radius is 8ft?
a.)
32 ft
25 ft
, b.)
,
25 min
32 min
c.)
75 ft
25 ft
20 ft
, d.)
, e.)
32 min
256 min
25 min
140
________________________________________________________________________
n =0
converge to a value of 3?
a.)
3
2
3
, b.) , c.) , d.) The series diverges, e.) Not enough information provided.
2
3
4
________________________________________________________________________
20.) The following graph represents the acceleration of an object:
Which of the following graphs could represent the position of the object?
141
a.)
b.)
c.)
d.)
e.)
21.) What is the average value of the function f ( x) = x 3 sin x on the interval [0, ] ?
a.) 3 6 1 , b.) 3 + 6 , c.) 2 6, d.) 3 6 , e.) 3 + 6
________________________________________________________________________
142
1 + 20 x 8 dx. If this
particular curve contains the point (1, 3), which of the following could be an equation for
this curve?
2 5 5 14
x +
5
5
20 9 7
y=
x +
9
9
2 5 5 15 2 5
y=
x +
5
5
2 5 5 15
y=
x +
5
2 5
4
y = 2x + 1
a.) y =
b.)
c.)
d.)
e.)
________________________________________________________________________
23.) A rectangle is to be inscribed between the curve y = 4 2 x 2 and the x-axis, with its
base on the x-axis. For which value of x will the area of this rectangle be maximized?
a.)
3
, b.)
2
2
, c.)
3
3
, d.)
2
2 , e.)
2
2
________________________________________________________________________
143
24.) Calculate the area between the curves y = e x and y = e 3 x , bounded by the lines y =
1 and y = 4.
8
10
a.) ln 4
3
3
8
b.) ln 4 4
3
8
c.) ln 4 2
3
1
2
d.) e12 e 4 + e
3
3
1 12
2
e.) e e 4 e
3
3
________________________________________________________________________
25.) What is the coefficient of x 3 in the Maclaurin series for
a.)
e2x
?
2
2
1
1
1
3
, b.) , c.) , d.) , e.)
3
6
2
3
6
________________________________________________________________________
144
________________________________________________________________________
27.) Consider the region bounded by the curve y = x 2 + 3, the line x = 1 , the x-axis, and
the y-axis. What is the volume of the solid generated when this region is revolved about
the y-axis?
a.) 88 , b.) 176 , c.)
24
3
17
, d.) , e.)
5
7
5
________________________________________________________________________
145
n =0
k n +1
, where k > 0?
( x + 3) n
________________________________________________________________________
End of Section I, Part A
If you finish before time is called, you may check your work in this section ONLY.
Do not proceed to Section I, Part B until told to do so.
146
Section 1, Part B
50 minutes
17 questions
Directions: Solve each of the problems in this section, using the available space for
scratchwork. Select the letter that corresponds to the best solution to the problem. No
credit will be awarded for anything written in the test booklet. Use your time wisely. A
graphing calculator is required for this section of the exam.
Note: 1.) Some of the numerical answers will not be exact. When this is the case, choose
the answer that best approximates the exact numerical value. 2.) Unless otherwise
indicated, the domain of a function f is the set of all real numbers x for which f (x) is a
real number.
29.)
x3 x dx =
2
3
ln 3
3x
ln 3 x 2
3x
+ C , c.)
3 , d.)
+ C , e.) 2 x3 x x + C
a.) 3 + C , b.)
2
2
2 ln 3
x
________________________________________________________________________
t5 +1
gallons per minute.
2t + 3t 2 + 2
How much water has passed through the turbine after one hour?
30.) The flow of water through a turbine is modeled as
147
en + 1
, what is lim u (n) ?
n
3n 1
1
, e.) ln(3 + e)
ln 3
________________________________________________________________________
32.) Consider the graph of the function f (x) below:
f ( x)dx exists
b
a.) I only, b.) III only, c.) I and II, d.) I and III, e.) I, II, and III
________________________________________________________________________
33.) Which of the following represents the value of the approximation y(1.2) of the curve
y = x ln x using a 4th-order Taylor polynomial centered at x = 1?
1631
409
1661
368
219
a.)
, b.)
, c.)
, d.)
, e.)
7500
1875
7500
1682
1000
________________________________________________________________________
148
dF
= 1.77t 2 + e 0.03t models how quickly the rate of a
dt
quantity changes. If the initial rate at which the quantity changes is 0, how much of the
quantity is there when t = 10 ?
34.) The differential equation
a.) 599, b.) 1187, c.) 589, d.) 1520, e.) 1177
________________________________________________________________________
35.) The approximate value of the area under the curve y = 2 x 2 2 x 3 + 2 x + 6 on the
interval [0, 1] can be found by using the trapezoidal rule. When using 5 trapezoids, what
is the approximate area under the curve on this interval using only the trapezoidal rule?
a.) 7.17, b.) 7.167, c.) 7.16, d.) 7.2, e.) 7.1667
________________________________________________________________________
5 x3
a.) 15 x sin(5 x )
b.) 15 x 2 cos(5 x 3 )
c.) cos(5 x 3 ) + cos(1)
d.) 5 x 3 sin(15 x 2 )
e.) 5 x 3 cos(15 x 2 )
149
________________________________________________________________________
d2y
37.) If the curve y(x) is defined parametrically as x = 3t 3 , y = 5 cos t ,
=
dx 2
a.) 5 sin t
45t 2 cos t + 90t cos t
b.)
729t 6
45t 2 cos t + 90t
c.)
81t 2
5 sin 2 t
d.)
9t 2
sin t cos t (9t 3 + 18t 2 )
e.)
9t 6
________________________________________________________________________
38.) Determine the area inside the circle r = 1 and outside the cardioid r = 3(1 + sin ) .
a.) 1.334, b.) 12.412, c.) 6.206, d.) 8.315, e.) 0.667
________________________________________________________________________
39.) Which of the following are -values at which there are horizontal tangents of
r = 4 sin on the interval [0, 2 ] ? (Note: Not all of the values may be listed.)
I.) = 0, II.) = , III.) 5.6642
a.) I only, b.) II only, c.) III only, d.) I, II, and III, e.) None of the above
150
________________________________________________________________________
40.) On the interval [0, 5], which of the following graphs could represent the function f
1 5
with the property that
f (t )dt = 2 ?
5 0
a.)
b.)
c.)
d.)
e.)
________________________________________________________________________
41.) Sand is poured on top of a conical pile at the rate of 200 ft 3 / sec , causing the radius
1
of the base of the pile to increase at ft/sec. How quickly is the height of the pile
2
increasing at the instant the radius of the base is 18 ft and the height of the pile is 10 ft?
151
________________________________________________________________________
42.) The rate of change of a quantity is given by the differential equation
dy
6t + 3
= 2
. If y (2) = 7, which of the following is the value of the constant term in
dt t + 2t 3
the equation for y?
a.) 7
15
9
15
9
15
ln 5 , b.) 7 ln 5 , c.) 7 ln 2 , d.) 7 ln 2 , e.) 7 ln 7
4
4
4
4
4
________________________________________________________________________
43.) If the velocity of a particle is given by the parametric equations
x(t ) = 3t 3 , y (t ) = sin 2 t , what is the speed of the particle at t = 1?
a.) 3.148, b.) 9.909, c.) 3.708, d.) 1.926, e.) 9.046
________________________________________________________________________
152
x A
a.) I only, b.) II only, c.) III only, d.) I and II, e.) II and III
________________________________________________________________________
45.) If y = 5 x 66 ,
d 66 y
=
dx 66
________________________________________________________________________
End of Section I, Part B
If you finish before time is called, you may check your work in this section ONLY.
Do not proceed to Section II, Part A until told to do so.
153
1.) Let the shaded region shown above be enclosed by the graphs of y = sin( x 2 ) and
y = x3 2x 2 .
a.) Calculate the area of the shaded region.
b.) Determine the volume of the solid generated when the shaded region is revolved
about the y-axis.
154
c.) Let the shaded region be the base of a solid. For this solid, each cross section
perpendicular to the x-axis is a semicircle. Calculate the volume of this solid.
________________________________________________________________________
2.) Two particle are moving in the xy plane. The velocity of Particle 1 is given
parametrically as:
v x = t 3 5t 2 3t
1
vy = e 2 .
The velocity of Particle II is given as v x = t 2 3, while v y is not explicitly given.
Both particles are initially at the origin.
a.) Determine x(t ) and y (t ) for Particle 1 and x(t ) for Particle 2.
155
b.) It has been found that y (t ) is an increasing linear function for Particle 2. If Particle 1
and Particle 2 have the same y-coordinate at t = 0.708399, find y (t ) for Particle 2.
d.) Do any of the particles change direction from t = 0 to t = 2 ? If so, find the total
distance traveled by the particle(s) when it (they) first change direction.
156
________________________________________________________________________
3.) Consider the differential equation
dy
= x + 2xy 2 .
dx
a.) Using Eulers method, approximate the value of y (1.4) by using increments of 0.1 if
y (1) = 0.5 .
b.) In the grid provided, draw the slope field that represents the family of solutions for
this differential equation.
157
y
3
x
3
c.) Rewrite in polar form the expression to which
dy
is equal.
dx
________________________________________________________________________
End of Section II, Part A
If you finish before time is called, you may check your work in this section ONLY.
158
4.) The graph above represents the rate at which the altitude of a traveling object changes.
Time t is on the horizontal axis and change in altitude A/t is on the vertical axis. The
graph is composed of line segments and a semicircle.
a.) On what intervals, if any, is the object gaining altitude? Justify your answer.
159
c.) On a set of axes (not provided), sketch the graph of the objects acceleration (A/t2)
versus time. The graph does not have to be to scale. Justify your sketch.
d.) Calculate the average rate of change of altitude of the object from t = 0 to t = 10.
160
________________________________________________________________________
5.) The following are true of the functions of f ( x) and g ( x) :
i.) f ( x) = x 2 g ( x)
f (1 + k ) f (1)
ii.) lim
=1
k 0
k
iii.) lim f ( x) = 2
x
f ( x)
= 1.
g ( x)
161
6.) A function f, which has a range of [-2, 1], is defined by the following power series for
all real numbers x:
n =0
(1) n 2 x n n
2x 4x 2 6x 3
=
+
+L.
(3n + 1)!
4!
7!
10!
a.) Find f (0) and f (0) . Does f (x) have a minimum, maximum, or point of inflection
at x = 0? Justify your answer.
162
________________________________________________________________________
End of exam
If you finish before time is called, you may check your work in this section and Section II
A, but you will not be allowed the use of a calculator.
163
3 x + y = 3 x 4 y 2
+ 4x3 y 3
dx
dx
dy
dy
3x + 3 y = 3x 4 y 2
+ 4x3 y 3
dx
dx
dy
4 2
(3 x 3 x y ) = 4 x 3 y 3 3 y
dx
dy 4 x 3 y 3 3 y
.
=
dx 3 x 3 x 4 y 2
2.) B. The carrying capacity is reached when the population stops growing. Since the
equation given represents the rate of change, one sets it equal to zero and solves for N.
N
The whole equation will equal zero when the term 1 equals zero.
72
3.) B. This is a definite integral in which one of the limits of integration is a variable:
4.) B. For a volume of solid of revolution problem, one has three possibilities: the disc
method, the washer method, or the shell method. The disc method would not work for
this problem because the radius would need to be perpendicular to the y-axis and, due to
the way in which the functions intersect, there would be no direct way to express the
radius algebraically. The washer method cannot be used because there is no hole in the
resulting solid. One must use the shell method. The formula associated with this method
when revolution is about the y-axis is V = 2 hr dx , where h is the height of the shell
a
(the distance between the bounding functions) and r is the radius of the shell (the distance
from the axis of rotation to the bounded functions; for a rotation about the y-axis, this
distance is simply x). Thus, the volume should be expressed as
2
5.5
(g ( x) f ( x) )
xdx .
5.) A. The question is asking for the slope of the tangent line, which is
dy
. For a set of
dx
dy
dy
3t 2
3
dy dt
=
= t 2 sec t.
parametric equations,
. Thus, in this instance,
=
dx 5 cos t 5
dx dx
dt
164
6.) E. This integral requires integration by partial fractions. Factoring the denominator,
one of the factors is linear, while the other is cubic, meaning that the numerators of the
partial fractions will be constant and quadratic, respectively:
Bx 2 + C
2x 3 + 2x 2 + 1 A
= + 3
x x + 2x 2 + 1
x 4 + 2x 3 + x
To solve for A, B, and C, one multiplies each side of the equation by the denominator:
2 x 3 + 2 x 2 + 1 = A( x 3 + 2 x 2 + 1) + Bx 3 + Cx
2 x 3 + 2 x 2 + 1 = Ax 3 + 2 Ax 2 + A + Bx 3 + Cx
2 x 3 + 2 x 2 + 1 = x 3 ( A + B) + x 2 (2 A) + x(C ) + A
Since the constant on the left is 1, this must also be the value of A, since it is the only
constant on the right. Since there are no terms with x on the left, this must also be the
case with the right, so C = 0. Finally, since the coefficient of x 3 on the left is 2, this must
also be the case on the right, meaning A+B = 2. Since A was already found to be 1, B
must also be 1. One can now integrate:
2x3 + 2x 2 + 1
dx =
x 4 + 2x3 + x
= ln x + ln x +
dx
+
x
x2
dx =
x3 + 2x 2 + 1
dx
+
x
x2
dx +
x3
x2
dx +
2x 2
x dx
2
1
1
1
1
x + x 3 + C = 2 ln x + x 3 + x + C.
2
3
3
2
7.) C. The best way to do this problem is to test each series separately. The first series
diverges by the nth term divergence test. Since the degree in the numerator and
denominator is the same, and since the limit is being taken at infinity, one takes the ratio
of the leading coefficients. Since this ratio is 1 in this case, one can conclude that the
series diverges because the limit is a number other than 0. The second series diverges by
2(n + 1)! 10 n
2(n + 1)n! 10 n
n +1
lim
the ratio test: lim
=
= lim
= . Since this limit is
n
n
+
+
(
1
)
(
1
)
n 10
2n! n 10
2n! n 10
greater than 1, the series diverges. The third series converges by the comparison test.
3
. Since r < 1, the geometric
4
n =1
series converges. Since this larger series converges, the original smaller series must also
converge. Thus, the third is the only convergent series.
One can compare it to the larger geometric series
8.) C. This problem just requires differentiation with the chain rule:
2
1
1
1
d 3
ln(3x) + 5 = [ln(3x) + 5] 3 =
.
3
dx
x 3 x 3 [ln(3x) + 5] 2
165
9.) C. This problem requires integration by parts. Since one term is a polynomial and the
other is an easily integrable function, one can use the tabular method:
u
3x 3
9x 2
18x
18
0
1
+1
-1
+1
-1
+1
-1
dv
cosx
sinx
-cosx
-sinx
cosx
From this, one finds that the integral is 3x 3 sin x + 9 x 2 cos x 18 x sin x 18 cos x + C .
10.) D. In order for the function to be continuous, the second part must begin where the
first part ended. The y-value at which the first part ended may be determined by
substituting 2 for x: (2) 2 1 = 3. Thus, the second part, which is a horizontal line, must
5
have a value of 3. Thus, one can find a: 3a 2 = 3 a = .
3
11.) E. While the graph of the derivative crosses the x-axis 6 times and is non-existant at
one point, this only means that the actual function has 7 critical points on the interval.
These critical points indicate relative extrema only. The function could also have
absolute extrema at the endpoints. Thus, a conclusion cannot be reached without more
information on the function.
12.) A. Attempting to substitute infinity for initially results in the indeterminate form
1
1
ln 2
1
1
1
0 0. Invoking a property of logarithms, lim 2 lim ln 2 = lim = .
1
ln 2
2
Thus, one uses LHpitals Rule: lim = lim = 0. This is not the answer,
however. Recall that the natural logarithm of the limit was taken. Thus, one must
reverse this by exponentiating: e 0 = 1. This is the final answer.
13.) C. In order to find the value that satifies the Mean Value Theorem, one must set the
slope of the tangent line (the derivative) equal to the slope of the secant line (the average
slope over the interval). For this problem,
y y (2) y (0) 13 1
msec =
=
=
= 6.
x
20
2
One then sets this value equal to the derivative to find at what value of x the derivative
equals this average slope:
166
dy
4
2
= 3x 2 + 2 = 6 x 2 = x =
. This is the value of c that satisfies the Mean
dx
3
3
Value Theorem.
14.) A. This integral can be manipulated to be in the form
du
, which equals
+ a2
1
u
tan 1 . This can be achieved by multiplying the integrand by 3 and multiplying the
a
a
1
outside by :
3
1
3
3dx
4 + 9x 2
= lim
a 1
11
3x
3x
3a 1
3
1
1
= lim tan 1 = lim tan 1 = lim tan 1
tan 1
2
a 3 2
a 6
2
2 1 a 6
2 6
2
4 + 9x
1
3dx
3
3 1
tan 1 .
=
2
2 12 6
3a 1
1
Thus, lim tan 1
tan 1
a 6
2
6
15.) B. In this problem, one is able to substitute x + 1 wherever there is an x in the series
and then multiply every term by x. Note that this exact procedure must be taken; one
cannot multiply every term by x and then substitute x + 1 for every x (This would actually
yield the answer represented by choice C). The reason for this is because one must
consider the manipulation of the functions argument first and only then consider what is
being multiplied. Thus, one must first consider the x in e x by replacing every x with
x + 1, and then multiply every term by x.
16.) E. The function g ( x) =
the area (including the sign negative or positive) under f (x), one can determine the
values of g (x). The areas associated with the graph are shown below:
Note that the left-most triangle has an associated area of -4 because the accumulation
function begins at x = 0. Thus, one must negate whatever area is to the left of that.
Starting at x = -2 and summing all of the areas up to x = 10, one would see that a value of
zero is never reached: 4 4 + / 2 / 2 + 6 = 2 . Thus, g (x) has no roots on the
interval.
167
derivative of g (x). Thus, the graph depicted represents the derivative of g (x). Since this
is the case, f (x) would be increasing on whichever interval g (x) is positive. The only
intervals on which f (x) is always positive are [-2, 0] and [3, 5]. Thus, f (x) is increasing
on these intervals.
18.) B. This is a related rates problem. Perhaps the best way to go about solving it is to
record all of the information given in the problem, write mathematically what the
problem is asking for, and write the equation to be used. In this problem, it is given that
dr
dV
ft 3
= 200
. One must find
when r = 8ft. The equation to be used is the
min
dt
dt
4
expression for the volume of a sphere: V = r 3 . Unfortunately, it is not guaranteed that
3
these sorts of geometric formulae will be given on the exam. One now differentiates the
equation, substitutes the known values, and solves for the unknown element:
d
d 4
dV
dr
dr
(V ) = r 3
= 4 r 2
(200 ft 3 / min) = 4 (8ft ) 2
dt
dt 3
dt
dt
dt
dr
256 ft 2
200ft 3 / min
dt dr = 25 ft .
=
2
2
dt 32 min
256 ft
256 ft
19.) A. This geometric series can be rewritten in the more familiar form
n =0
1
k .
2
1
< 1, the series certainly converges. This is a geometric series, so its sum has
2
a
. Since S n and r are known, one can find a:
the formula S n =
1 r
a
3
(3) =
a= .
2
1
1
2
Since r =
20.) D. This is a very concept-based question. It concerns the relationship between the
graph of a function (position) and its second derivative (acceleration). For this kinds of
problems, it is best to observe each part of the graph separately and conclude such
relationships. Beginning with the first part of the acceleration graph, the acceleration
linearly decreases until it reaches zero. On the velocity graph, one would see a quadratic
increase with linearly decreasing slope until a relative maximum is reached. On the
position graph, one would see a cubic increase until a point of inflection is reached.
Thus, one can rule out choices E and B because they do not exhibit a point of inflection at
the value where the acceleration is zero. Analyzing the acceleration graph further, the
168
acceleration is negative and linearly decreasing until it becomes constant. After this
constant value, the acceleration begins to increase toward zero. One the velocity graph,
one would see a quadratic decrease in velocity with linearly decreasing slope and then a
linear decrease in velocity, followed by a quadratic decrease in velocity with linearly
increasing slope. On the position graph, one would see a cubic increase in position with
quadratically decreasing slope until another point of inflection is reached. This leaves
only choice D.
21.) C. Recall that the average value of a function over a certain interval [a, b] has the
b
f ( x)dx
ba
integration by parts. However, since one term is a polynomial and the other is an easily
integrable function, one can use the tabular method:
u
x3
3x 2
6x
6
0
dv
sinx
-cosx
-sinx
cosx
sinx
1
+1
-1
+1
-1
+1
-1
From this,
x sin xdx = x cos x + 3x sin x + 6x cos x 6 sin x + C = cos x(6x x ) + sin x(3x
3
6) + C.
x 3 sin xdx
[cos( )(6
) + 0 [cos(0)]
= 2 6.
22.) C. Recall that the arc length of a curve has the formula l =
b
a
dy
1 + dx. With
dx
dy
this formula in mind, one can see that = 20 x 8 for this problem. Thus, this
dx
problem is essentially asking one to solve a differential equation and use a set of initial
conditions:
dy
2 5 5
= 20 x 4 = 2 5 x 4 y = 2 5 x 4 dx =
x + C.
dx
5
169
2 5 5
15 2 5
(1) + C = 3 C =
.
5
5
2 5 5 15 2 5
.
Thus, y =
x +
5
5
It is given that y (1) = 3, so
23.) B. This is an optimization problem, which involves setting the derivative of a certain
equation equal to zero and finding the values that make the original function the largest
or smallest that it can be on a certain interval. In this problem, one wishes to find the
value that makes the area function have the greatest magnitude. To do this, one
determines the derivative of the area function, sets this expression equal to zero, and finds
the value at which the derivative shifts from being positive to being negative (i.e. where
the area function has a relative maximum). The area in question is that of a rectangle
with its height equal to the distance from the x-axis to the function, which is merely
y = 4 2 x 2 , and its length equal to 2x, since the base is symmetric about the y-axis.
Thus, A = lh = (2 x)(4 2 x 2 ) = 8 x 4 x 3
dA
2
= 8 12 x 2 = 0 x = . Since at
dx
3
2
the derivative shifts from positive to negative values, the area function has a
3
relative maximum at this point.
x=+
24.) C. Since no calculator may be used on this portion of the exam, one must use a little
mathematical intuition here. Since y = e 3 x increases more quickly than y = e x , the
former curve will generally lie closer to the y-axis than the latter curve. Thus, y = e 3 x is
farther to the left and y = e x is farther to the right. Since the boundaries in this case are
horizontal lines, there is no top and bottom function that one can use to determine the
area. However, there is a left and right curve, but one must have everything,
including the limits of integration, in terms of y since the infinitesimal rectangles are now
perpendicular to the y-axis instead of the x-axis. The leftmost curve in terms of y is
ln y
x=
and the rightmost curve in terms of y is x = ln y. One now evaluates the definite
3
4
4
4
ln y
1 4
integral: A =
" Right"" Left" dy =
ln
y
dy
=
ln
y
dy
ln y dy. To
3
3 1
1
1
1
integrate the natural logarithmic function, one technically needs to use integration by
parts. However, I recommend memorizing what the integral is to save time:
8
1
4
4
ln y dy = y ln y y + C . Thus, A = ( y ln y y ) 1 ( y ln y ) 1 = ln 4 2.
3
3
25.) A. This question concerns the algebraic manipulation of power series. Note that the
Maclaurin series for e x is one of the power series that should be memorized for the exam.
The others can be found on page 61. For this problem, one must substitute 2 x for x and
170
1
multiply by . Since the Maclaurin series for e x is
2
the Maclaurin series for
1 2x
e is
2
n =0
n =0
xn
x2 x3 x4
= 1+ x +
+
+
+ L,
n!
2! 3! 4!
(2 x)
1
4x 2
8 x 3 16 x 4
= +x+
+
+
+ L. Thus, the
2n!
2
2 2! 2 3! 2 4!
n
8
2
coefficient of x 3 is
= .
2 3! 3
26.) B. One must first differentiate this expression implicitly:
d
d
dy
dy
dy
dy
(2 xy) =
( x + y ) 2 x + y = 1 +
2x + 2 y = 1 +
dx
dx
dx
dx
dx
dx
dy
dy 1 2 y
(2 x 1) = 1 2 y
.
=
dx
dx 2 x 1
dy
in polar form, one must recall that x = r cos and y = r sin . So,
dx
dy 1 2 y 1 2(r sin ) 1 2r sin
=
=
=
.
dx 2 x 1 2(r cos ) 1 2r cos 1
To express
the y-axis is V = 2 hr dx , where h is the height if the shell (the distance between the xa
axis and the function) and r is the radius of the shell (the distance from the axis of
rotation to the bounded function; for a rotation about the y-axis, this distance is simply x).
Thus, the formula for the volume of this solid is
1
17
3
1
.
V = 2 ( x + 3) xdx = 2 x 5 + x 2 =
5
2 0
5
0
28.) B. Determining the radius of convergence of a power series requires using the ratio
test and determining for which values of x the series converges. Perhaps the only more
slightly complicated aspect of this problem is that the constant k is unknown.
171
n =0
k n+ 2
n +1
k n +1
( x + 3)
lim
n
( x + 3) n
k n +1
n
( x + 3)
k n + 2 ( x + 3) n
= lim
n ( x + 3) n +1
k n +1
k
=
.
x+3
a n +1
< 1. Thus,
n a
n
k > x + 3 > k ( k 3) > x > (k 3). This represents the radius of convergence.
Section I B:
29.) D. The integral of a general exponential a u is given as
a du = ln a a
u
+C.
For x3 x dx, one must first multiply the integrand by 2 to give du = 2 xdx and
multiply the outside of the integral by
1
to compensate for this. The solution to the
2
1 1 x2 3 x
3 =
.
integral is, thus:
2 ln 3
2 ln 3
30.) B. Recall that the definite integral of a rate yields a total amount. Since the equation
in this problem represents flow, integration over a certain time period will yield the total
amount associated with that time period. One should perform the following definite
60
t5 +1
integral on the calculator:
dt = 1063878.5 gallons.
2
0 2t + 3t + 2
31.) A. No tests are really required for this problem. In both the numerator and
denominator of this fraction are exponentials. However, the denominator contains an
exponential whose base (3) is larger than that in the numerator ( e 2.71828 ). Thus, the
denominator will increase more quickly than the numerator; the addition or subtraction of
1 is negligible. The effect of the more rapid increase in the denominator is that the
fraction approaches zero as n becomes larger and larger.
32.) C. Only the first two statements are correct. The first statement is correct because
on the open interval (a, b), the function is decreasing, meaning the derivative is negative
on that interval. The second statement is correct because at a, the function exhibits a
corner, meaning the function is not differentiable at that point. The third statement is not
necessarily true. While it does seem that the graph will continue to decrease as x
becomes larger and larger, implying that the area under the curve converges on a finite
172
value, this is not guaranteed. Without more knowledge of the functions behavior, one
cannot truly know whether the function continues to decrease toward a y-value of 0.
33.) C. One must find the first several terms of the polynomial expansion for y = x ln x
to approximate y (1.2). Recall that the formula for a Taylor expansion is
n =0
f ( n ) (c)( x c) n
. It is given that c = 1. Applying the formula to y = x ln x up to n = 4,
n!
2!
3!
4!
2!
3!
4!
1661
0.2 + 0.02 0.001333 + 0.0001333
.
7500
34.) D. This question requires careful reading. It states that the equation given models
how quickly the rate of a quantity changes. Thus, it is actually measures the rate of rate!
This means that the equation given represents the second derivative of the equation
representing the amount, meaning one must integrate twice. Integrating the first time
yields the equation representing the rate of change of the quantity, or the first derivative
of the equation representing the amount. It is stated that the initial rate is zero. With
these initial conditions, one may find the exact equation representing the first derivative:
First derivative =
(1.77t
To determine the amount at t = 10, one integrates the equation for the first derivative
from t = 0 to t = 10. Using the graphing calculator:
10
173
A =
1
2
d 5x
f ( x) =
sin tdt = 15 x 2 sin(5 x 3 ) .
dx 1
37.) B. Recall that the formula for the second derivative of a parametrized curve is:
d dy
2
dy
d y dt dx
=
. One first determines
:
2
dx
dx
dx
dt
dy
dy dt 5 sin t
. One then differentiates this expression with respect to t:
=
=
dx dx
9t 2
dt
d dy (9t 2 )(5 cos t ) (5 sin t )(18t ) 90t sin t 45t 2 cos t
=
=
dt dx
81t 4
81t 4
dx
:
dt
174
To determine where the curves intersect, one sets them equal and solves for :
2
= 0.72972766. This angle is measured in a
3
clockwise sense. To express it in a counterclockwise sense, one adds it to 2 (i.e.360 o ) .
Thus, one angle is 2 0.72972766 = 5.553458. Since the value of sine is negative in
quadrants III and IV, the other angle must be in quadrant III, 0.72972766 radians from
(180 o ) . Thus, the second angle is + 0.72972766 = 3.871320. Now that the two
limits of integration have been found, one can compute the area:
1 = 3 + 3 sin sin =
1 2
1 5.553458
1
r2 r12 d =
(3 + 3 sin ) 2 1 d = (1.333963) 0.6669815
2
2 3.871320
2
(Note that though the integral is negative, area must always be positive).
A=
dy
dx
= 0 and where
0 at the
d
d
same time. One first determines the derivative of y with respect to :
signs, since the average value concerns the integral) and divides this value by 5. The
only graph whose total area divided by 5 is 2 is the one depicted in choice E.
41.) C. This is a related rates problem. As always, one should record every quantity
given, determine what must be found, and find the equation to be used. In this question,
dV
dr 1
= 200ft 3 / sec,
= ft/sec, and that when r = 18ft, h = 10ft. The
it is given that
dt
dt 2
dh
at these values. For this problem, one can use the
question is asking for the value of
dt
1
equation for the volume of a cone: V = r 2 h . Do not expect this formula to be given
3
on the exam; it is better to memorize it and those for other geometric solids. One
differentiates this equation, substitutes all necessary values, and solves for the unknown:
d
d 1
dV 1 2 dh
dr
(V ) = r 2 h
= (r ) + (h) 2r Note that the product rule was used here.
dt
dt 3
dt 3 dt
dt
1
2 dh
(200ft 3 / min) = (18ft ) + (10ft)(2(18ft)(1/2ft/min) )
3
dt
3
2 dh
3
(200ft / min) = 324ft
+ 180ft / min
3
dt
600 3
2 dh
3
ft / min = 324ft
+ 180ft / min
dt
dh
= 0.0339ft/min.
dt
42.) A. Integration by partial fractions is required to solve this differential equation:
6t + 3
6t + 3
A
B
dt 2
=
+
t + 2t 3
t + 2t 3 (t + 3) (t 1)
6t + 3 = A(t 1) + B(t + 3) = At A + Bt + 3B = t ( A + B) + 3B A
y=
Since the constant on the left side is 3, 3B A = 3 . Since the coefficient of t on the left
side is 6, A + B = 6. Solving these two equations simultaneously, one finds that
15
9
A = and B = . One can now integrate:
4
4
y=
6t + 3
15
dt =
4
t + 2t 3
2
dt
9
+
t +3 4
B
15
9
= ln t + 3 + ln t 1 + C
t 1 4
4
15
9
15
ln (2) + 3 + ln (2) 1 + C = 7 C = 7 ln 5.
4
4
4
176
dx dy
43.) E. The speed ( v ) of a particle is given parametrically as v = + .
dt dt
Thus, for this particle,
dx
dy
= 9t 2 ,
= 2 sin t cos t.
dt
dt
2
v = 9t 2 + (2 sin t cos t ) 2 = 81t 4 + 4 sin 2 t cos 2 t At t = 1,
v = 81 + 4 sin 2 (1) cos 2 (1) 9.045818.
( )
44.) E. The first statement is not necessarily correct because a definite integral may very
well exist over an interval on which there are discontinuities. The second statement is
true; approaching A from the left side or the right side yields the same limit. The third
statement is also true; since there is a removable discontinuity at A, the derivative does
not exist there. Remember that continuity is a pre-requisite for differentiability.
45.) B. Obviously it would be quite impractical to evaluate 66 derivatives! Instead, one
should evaluate the first several derivatives and find a pattern that can be used to find the
66th derivative:
dy
= 5(66) x 65
dx
d2y
= 5(66 65) x 64
2
dx
d3y
= 5(66 65 64) x 63
3
dx
d4y
= 5(66 65 64 63) x 62
4
dx
Note that the order of the derivative corresponds to how many new terms are
being multiplied. Thus, when the 66th derivative is taken, there will be 66 multiplied
d 66 y
= 5(66!) .
terms. This means that
dx 66
Section II A:
1.) a.)
A=
x2
("Top"-"Bottom")dx
x1
1.8862947
[(sin( x
)) ( x 3 2 x 2 ) dx = 2.1810975 1.181.
177
x2
1.8862947
c.) V =
x1
x2
1 2
1 (sin( x 2 )) ( x 3 2 x 2 )
Semicircle A = r The radius is 1/2 the base A =
2
2
2
1
V =
8
1.8862947
[(sin( x
2
1
)) ( x 3 2 x 2 ) dx = (2.181097499) = 0.856514985 0.857
8
Rubric:
Part A: 1 point for setting up integral correctly
1 point for correct limits of integration
1 point for correct answer to 3 decimal places
Part B: 1 point for specifying shell method
1 point for setting up integral correctly
1 point for correct answer to 3 decimal places
Part C: 1 point for volume equation
1 point for setting up integral correctly
1 point for correct answer to 3 decimal places
Total: 9 points
2.) a.) Particle 1:
dx
dt
1
5
3
3
2
x = (t 5t 3t )dt = t 4 t 3 t 2 + C
4
3
2
1 4 5 3 3 2
1
5
3
(0) (0) (0) + C = 0 C = 0 x(t ) = t 4 t 3 t 2
4
3
2
4
3
2
v x = t 3 5t 2 3t =
1
t
2
dy
dt
1
t
12 t
y = e dt = 2e 2 + C
vy = e
178
2e 2
(0)
+ C = 0 C = 2 y (t ) = 2e 2 2
Particle 2:
dx
dt
1
2
(t 3)dt = t 3 3t + C
3
vx = t 2 3 =
x=
1 3
1
(0) 3(0) + C = 0 C = 0 x(t ) = t 3 3t
3
3
b.) If Particle 2s y-position is an increasing linear function, let y be generally represented
as y (t ) = at + b . Since y (0) = 0, b = 0 .
y-coordinate of Particle 1 at t = 0.708399:
y (0.708399) = 2e
1
( 0.708399 )
2
2 = 0.85007891.
c.) Particle 1 : a x (t ) =
dv x
= 2t
dt
dv y
a y (t ) =
=0
dt
Particle 2: a x (t ) =
a x (1.5) = 2(1.5) = 3
a(t ) = 3i.
179
d.) Yes, Particle 2 changes direction at t = 1.7320508 because the x-velocity is zero at
that time.
Total distance =
Total distance =
b
a
1.7320508
0
dy
dx
+ dt
dt
dt
6
t 3 + dt = 4.1592971 4.159.
5
2
Rubric
Part A: 1 point for correctly setting up the differential equation for each particle
1 point for correctly solving the differential equation for each particle
Part B: 1 point for providing the skeleton equation for y, or specifying that one must
solve for an unknown constant
1 point for correctly solving for the unknown constant (either in decimal form or fraction
form)
1 point for providing the correct equation for y
Part C: 1 point for setting up the equation to solve for the acceleration vector for each
particle
1 point for providing the correct acceleration vector for each particle
Part D: 1 point for stating that Particle 2 changes direction at t = 1.7320508
1 point for correctly calculating the total distance traveled to 3 decimal places
3.) a.)
Total= 9 points
x1 = x 0 + x = (1) + (0.1) = 1.1
y 0 = (1) + 2(1)(0.5) 2 = 1.5
y1 = y 0 + xy 0 = (0.5) + (0.1)(1.5) = 0.65
x 2 = x1 + x = (1.1) + (0.1) = 1.2
y1 = (1.1) + 2(1.1)(0.65) 2 = 2.0295
y 2 = y1 + xy1 = (0.65) + (0.1)(2.0295) = 0.85295
x3 = x 2 + x = (1.2) + (0.1) = 1.3
y 2 = (1.2) + 2(1.2)(0.85295) 2 = 2.946056886
y 3 = y 2 + xy 2 = (0.85295) + (0.1)(2.946056886) = 1.147555689
180
y (1.4) 1.195
b.) While the tables below are not necessary to receive credit, one must show some sort
of calculations that led to the slope field.
dy/dx
(0,0)
0
(0,1)
0
(0,2)
0
(0,3)
0
dy/dx
(1,0)
1
(1,1)
3
(1,2)
9
(1,3)
19
dy/dx
(2,0)
2
(2,1)
6
(2,2)
18
(2,3)
38
dy/dx
(3,0)
3
(3,1)
9
(3,2)
27
(3,3)
57
The slope field should look approximately like this. One is not graded on hand-eye
coordination! As long as the slope field shows a sufficient basis in calculations, full
credit should be received.
c.)
dy
= x + 2xy 2
dx
Polar form: x = r cos , y = r sin
dy
= (r cos ) + 2(r cos )(r sin ) 2
dx
dy
= r cos + 2r 3 sin 2 cos .
dx
Rubric
Part A: 1 point for using the correct formulae for Eulers method
2 points for substituting the correct numerical values for the variables in these equations
1 point for providing the correct approximation to 3 decimal places
181
Note that the 4 came from subtracting the area of the trapezoid from the area of the
neglected triangle that makes up that trapezoid. Summing these areas,
Atotal
0,5
with an integral here (because the definite integral of a rate over a certain time period
yields a total amount) , so the signs of the areas (i.e. negative and positive) must be
preserved.
c.) This part required a lot of graphical and analytical intuition. In the end, the sketch
should look something like this:
182
Note that there are no tick marks because the problem specified that the sketch need not
be to scale.
Justification: The graph of A/t initially increases linearly, meaning A/t2 (the
derivative) has some constant positive value. A corner then appears, implying a point of
non-differentiability. After this point, A/t decreases linearly, meaning the A/t2 is some
constant negative value. The graph of A/t experiences another corner, meaning it has no
derivative at that point. The graph of A/t then becomes constant for a period of time,
meaning A/t2 is zero for that same time period. After another corner appears, the graph of
the A/t increases linearly, meaning A/t2 is a constant positive value. After another corner
appears, the slope tangent to the semicircle of A/t decreases. Since the semicircle is a
square-root function and the derivative of a square-root function is the reciprocal of a
square-root function, the graph of A/t2 decreases in a reciprocal-square root fashion.
When the graph of A/t reaches a maximum value at the top of the semicircle, the
reciprocal square-root function of A/t2 crosses the t-axis, and continues to decrease. After
one final corner appears, the A/t has a constant value, meaning A/t2 is zero.
b
f ( x)dx
ba
over the entire interval, one must sum all of the areas (including signs, since one is
dealing with an integral) under the curve and divide by the magnitude of the interval.
While the area of the triangles, trapezoids, and rectangles are simple to calculate, the
semicircular region is a bit troublesome. The area under the curve is not the area of the
semicircle because, since the circular curve is convex rather than concave (i.e. forming a
bowl), the actual area of the semicircle lies below the curve. Thus, one must find the area
of the imaginary rectangle in which the semicircle is inscribed and subtract from this area
the area of the semicircle. After doing all of the geometry, one can sum the areas
(including signs) and divide this by 10 to determine the average value of A/t:
183
Atotal
0 ,10
= 9528+
4=
10.
10
8
= 1.
Average value of A/t =
10
80
Rubric
Part A: 1 point for stating that the object is gaining altitude on [0, 3]
1 point for giving the correct justification
Part B: 1/2 point for summing the areas correctly
1/2 point for supplying the correct total area (total accumulated altitude)
Part C:
2
point for every correctly sketched portion of A/t2 (There are five portions)
5
2
point for every correct justification corresponding with each portion of the graph
5
(There should be a justification for every portion of the graph)
Part D: 1 point for indicating the integral equation for average value
1 point for correctly calculating the average value
Total: 9 points
5.) This is a very conceptual question. This may be good thing for some and signal utter
doom for others (like me the first time I encountered such a question on a test!). It
requires a painfully close analysis of each statement and how the statements relate to each
other. In essence, here is what they mean:
i.) The derivative of f (x) is equal to x 2 g ( x) (simple enough!)
ii.) The derivative of f (x) evaluated at x = 1 is 1 (This is the limit definition of the
derivative)
iii.) As x approaches infinity, f (x) approaches 2.
iv.) f (x) evaluated at 1 is 0 and the second derivative of f (x) evaluated at 1 is 1.
It is now time to synthesize all of this information in each part of the problem:
a.) Since f ( x) = x 2 g ( x),
integral,
(x
(x
g ( x) dx = lim
(x
a
a 1
Statement iii, states that lim f ( x) = 2 and statement iv states that f (1) = 0. Thus,
x
184
f ( x) f (1) 0
f ( x)
f ( x)
=
= Use LHpitals Rule: lim
= lim
x
1
x
1
g ( x) g (1) 0
g ( x)
g ( x)
f ( x) = x 2 g ( x) f ( x) = 2 x g ( x) g ( x) = 2 x f ( x)
g (1) = 2(1) f (1) = 1
lim
x 1
f ( x) f (1) 1
=
= =1
g ( x) g (1) 1
Rubric
Part A: 1 point for showing that
(x
g ( x) dx = f ( x) + C
d
6.) a.)
dx
(1) n 2 x n n
2 8 x 18 x 2
= +
+L
(3n + 1)!
4! 7! 10!
d
dx
(1) n 2(0) n n
2
=
(3n + 1)!
4!
n =0
n =0
d2
dx 2
d
dx 2
n =0
n =0
(1) n 2 x n n 8 36 x
=
+L
(3n + 1)!
7! 10!
(1) n 2(0) n n 8
=
(3n + 1)!
7!
185
Since neither the first derivative nor the second derivative is zero when x = 0, there
are no relative extrema or points of inflection at this point.
b.) Max error
1
100
1
f ( n +1) ( )( x c) n +1
. With a range of [-2, 1], f ( n +1) ( ) has the
Rn ( x ) =
100
(n + 1)!
greatest magnitude when it is 1. Thus,
1
(1)(1 0) n +1
1
.
=
100
(n + 1)!
(n + 1)!
1
1
1
1
=
and
=
, which is too large.
(5 + 1)! 120
(4 + 1)! 24
Therefore, one must use a mimimum of a 5th-order Maclaurin polynomial.
c.) Substitute every x with x 3 + 2 :
n =0
(1) n 2( x 3 + 2) n n
4th
(3n + 1)!
order f ( x 3 + 2)
2( x 3 + 2) 4( x 3 + 2) 2 6( x 3 + 2) 3 8( x 3 + 2) 4
+
+
4!
7!
10!
13!
Rubric
186
187
Step 5: Add up the total number of points achieved for the whole
free-response section.
Step 6: Multiply the total multiple-choice score by 1.2. This is the
multiple-choice composite score.
Step 7: Multiply the total free-response score by 1.0. This is the
free-response composite score.
Step 8: Add the results of steps 6 and 7.
Step 9: Find the AP score on the table below.
Composite Score Range
63-108
53-62
42-52
25-41
0-24
AP Grade
5
4
3
2
1
Note that this is my formula based upon how difficult I perceived the pratice exam in
comparison to past exams. This formula will almost certainly differ from the one
actually used.
188
( )
log10 (a) = log(a) (When a subscript is not written, the common logarithm is implied.)
log x ( x) = 1 (e.g. ln(e) = 1) (This is the case because the logarithmic function is the inverse
of the exponential function, and vice versa.)
x log x ( a ) = a (For the same reason as above)
Laws of Exponents
(a x )(a y ) = a x + y
(a x )
= a x y
y
(a )
1
a x = x (a 0)
a
Aspects of Square-Root Notation:
n
m
n
a , where a is the radicand, n is the index, and m is the power of the radicand.
Trigonometric Indentities:
sin 2 + cos 2 = 1
tan 2 + 1 = sec 2 (Pythagorean identities)
cot 2 + 1 = csc 2
cos 2 = cos 2 sin 2
sin 2 = 2 sin cos
1 + cos 2
(Double-angle formulae)
cos 2 =
2
1 cos 2
sin 2 =
2
Values of Trigonometric Functions to be Memorized:
189
sin(0) = 0
cos(0) = 1
tan(0) = 0
1
sin =
6 2
3
cos =
2
6
1
tan =
3
6
2
sin =
2
4
2
cos =
2
4
tan = 1
4
sin =
3
cos =
3
tan =
3
3
2
1
2
3
sin = 1
2
cos = 0
2
tan =
2
190
du = u + C
adu = au + C
du
u 1
2
= sec 1 u + C
(du dv) = du dv = u v + C
u n +1
+ C , (n 1)
n +1
u n du =
du
= ln u + C
u
a
a du =
, a > 0 and a 1
ln a
e u du = e u + C
u
cos
udu
= sin ud
References
Braden, Bart. Calculating Sums of Infinite Series. The American Mathematical
Monthly, 99(7), Aug. Sep. 1992: 649-655.
Bronson, Richard. Schaums Outline of Theory and Problems of Differential Equations:
Second Edition. New York: McGraw-Hill, 2003
Campbell, Neil A and Jane B. Reece. Biology: Seventh Edition. San Francisco:
Pearson: Benjamin Cummings, 2005.
Chihara, Charles S. On the Possibility of Completing an Infinite Process. The
Philosophical Review, 74(1), Jan 1965: 74-87.
Gray, Alfred. Modern Differential Geometry of Curves and Surfaces with
Mathematica: Second Edition. Boca Raton, FL: CRC Press, 1998.
Kahn, David S. Cracking the AP Calculus AB & BC Exams: 2002-2003 Edition.
New York: Princeton Review Publishing, 2002.
King, Kerry J. and Dale W. Johnson. Cliffs AP Calculus AB and BC: 3rd Edition.
New York: Hungry Minds, 2001.
Metz, Clyde R. Schaums Outline of Theory and Problems of Physical Chemistry:
Second Edition. New York: McGraw-Hill, 1989.
Pelcovits, Robert A. How to Prepare for the AP Physics C Examination. Hauppauge,
NY: Barrons Educational Series, 2002.
Peleg, Yoav; Reuven Pnini; and Elyahu Zaarur. Schuams Outline of Theory and
Problems of Quantum Mechanics. New York: McGraw-Hill, 1998.
Rubinow, S.I. Introduction to Mathematical Biology. New York: Dover, 2002.
Silberberg, Martin S. Chemistry: The Molecular Nature of Matter and Change: Fourth
Edition. New York: McGraw-Hill, 2006.
Tinker, Michael and Robert Lambourne. Further Mathematics for the Physical Sciences.
Chichester, UK: John Wiley & Sons, 2000.
Weir, Maurice D.; Joel Hass; and Frank R. Giordano. Thomas Calculus: Eleventh
Edition. Boston: Pearson: Addison-Wesley, 2005.
Wrede, Robert C. and Murray Spiegel. Schaums Outline of Theory and Problems of
Advanced Calculus: Second Edition. New York: McGraw-Hill, 2002.
192
Index
Absolute convergence, 42
Agnesi, Maria Gaetana, 82
Alternating series, 46, 52-53
Analytic geometry, 75
An Essay on the Principle of Population, 31
Archimedean spiral (see Spirals)
Arc length (in Cartesian coordinates), 75-78
of parametric curves, 84
of polar curves, 100
Astroid, 81
AP Calculus BC, 1
Bernoulli, Jakob, 91
Bounded growth (see Learning curve)
brachistochrone
Carbon-14, 34
Cardioid, 91
Carrying capacity, 28
Cartesian coordinates, 25
Cauchy, Augustin Louis, 66
Cauchy form of the remainder, 66
Cauchys Mean Value Theorem, 8
Center of mass, 7
Circle
parametric, 80
polar, 88
Closed interval, 19
Closed path, 129
Combinatorics, 45
Comparison test (see Infinite series)
Computational cost, 71
Conditional convergence, 42
Conservative field, 130
Conservative force (see Conservative field)
Convergent
improper integrals, 19
infinite series, 40, 42
Cross product (see Vector product)
Curl, 128
Curvature, 123
Del operator, 127
Descartes, Ren, 75
Kinetic-molecular theory, 22
Parallelogram method, 106
L'Analyse des Infiniment Petits pour
Parameter, 78
l'Intelligence des Lignes Courbes, 8
Parametric equations and curves, 79-87
Lagrange, Joseph Louis, 66
differentiation of, 82-84
Lagrange form of the remainder, 66
integration of, 84-87
Law of exponential change, 28
Partial derivative, 71, 127
Law of sines and cosines, 106
Partial fraction decomposition, 15
Learning curve, 29
Partial sum, 41
Leibniz, Gottfried, 6
Path independence, 130
Lemniscate, 91-92
Path integral (see Line integral)
Lennard-Jones, John, 69
Pi (case study of), 54
Lennard-Jones potential, 69
Pole, 87
LHpital, Guillaume de, 8
Polar equations and curves, 87-102
LHpitals Rule, 8-11
differentiation of, 94-98
Libby, Willard, 34
integration of, 96-102
Limaon, 90-91
Population ecology, 30
Power series, 56-74
Limit comparison test (see Infinite series)
Linear algebra, 17, 125
functions defined by, 58
Line integral, 128-130
interval of convergence of, 57
Linear restoring force, 68
method of (differential
LIPET, 12
equations), 72
Logistic growth, 28
radius of convergence of, 56
Principal unit normal vector, 123
Logarithmic spiral (see Spirals)
Maclaurin, Colin, 62
Probability density function, 22
Maclaurin polynomial (see Maclaurin series) Products (chemical), 33
Maclaurin series, 62
P-series, 44
194
Quantum mechanics, 23
Radial acceleration, 119
Radioactive decay, 34
Radiometric dating, 34
Radius of convergence (see Power series)
Ramanujan, Srinivasa, 54
Random variable, 22
Rate law, 34
Ratio test, 57 (see also Power series)
Reactants (chemical), 33
Reaction kinetics, 33
Reaction order, 34
Rectangular coordinates (see Cartesian coordinates)
Reduced row echelon form (RREF), 17
Remainder, 65
Repeated factor, 16
Resultant, 104
Right-handed coordinate frame, 111-112
Right hand rule, 114
Rose curve, 92-93
Scalar field, 126-128
Scalar field function (see Field function)
Scalar product, 112-114
Separation of variables, method of, 25
Series (see Infinite series)
Sequences (see Infinite sequences)
Simple harmonic motion, 68-69, 74
Speed, 119
Spirals, 89-90
Step size, 25
Stoichiometric coefficient, 33
Surface area of revolution
of parametric curves, 86
of polar curves, 101
Tangential acceleration, 119
Tautochrone (see Brachistochrone)
Taylor, Brook, 11, 63
Taylor polynomial, 64
Taylor series, 63-74
applications of, 68-74
Taylors theorem, 65
Telescoping series, 44
Terminal point (parametric), 79
Terminal velocity, 39
Term-by-term
differentiation, 59
Term-by-term
integration, 59
TI-89 calculator, 3
TNB frame, 125
Torricelli, Evangelista, 21
Torque, 115-117
Torsion, 125
Transcendental function,
10, 31
Treatise on Fluxions, 62
Truncation, 51-52
Truncation error
(see Truncation)
195
196
197
198
199