Phan Tich Cong C Dai Hocenglish

Bùi Xuân Vững Instrumental Analysis Chemistry
PART A: SPECTROCHEMICAL METHODS OF ANALYSIS

AN INTRODUCTION TO SPECTROSCOPY
Historically, the term ‘spectroscopy’ refered to a branch of science in which light
was resolved into its component wavelenghths to produce spectra, which were plots of
some function of radiant intensity versus wavelength or frequency. Nowsaday, the
meaning of spectroscopy has been broadened to include studies not only with visible
radiation but also other types of electroradiation, such as X ray, ultraviolet, infrared,
microway and radio frequency radiation.
Spectroscopy has played a vital role in the development of modern atomic
theory. In addition, spectrochemical methods have provided perhaps the most widely
used tools for the elucidation of the structure of molecular species as well as the
qualitative and quantitative determination of both inorganic and organic compounds.
I. General properties of electromagnetic radiaton
Electromagnetic radiation is a type of energy that takes numerous forms, the
most easily recognizable being visible light and radiant heat. The others include
gamma ray, X ray, ultraviolet, microwave, and radio-frequency radiation. Many of
properties of electromagnetic radiation are conveniently described by means of a
classical sinusoidal wave model, which employs such parameters as wavelength,
frequency, velocity, and amplitude. In contrast to other wave phenomena, such as
sound, electromagnetic radiation requires no supporting medium for its transmission
and thus passes readily through a vacuum.
The wave model fails to account for phenomena associated with the absorption
and emission of radiant energy. To understand these processes, it is necessary to
invoke a particle model in which electromagnetic radiation is viewed as a stream of
discrete particles called photon with the energy of a photon being proportional to the
frequency of the radiant. These dual views of radiation as particles and as waves are
not mutually complementary. Indeed, the duality is found to apply to the behavior of
streams of electrons and other elementary particle such as protons and it is completely
rationalized by wave mechanics.
1. Wave Properties of Electromagnetic Radiation
For many purposes, electromagnetic radiation is conveniently represented as
electric and magnetic field that undergo in-phase, sinusoidal oscillations as can be
seen in Figure .
Here only the electric component of radiation will be considered because the
electric field is responsible for most of the phenomena that are of interest to us,
including transmission, refraction, and absorption. Is is noteworthly, however, that the
1
magnetic component of electromagnetic radiation is responsible for the absorption of

radio-frequency waves in nuclear magnetic resonance.
Figure 0.1: (a) The in phase sinusoidal oscillations and (b) the electric component of radiation
2. Wave Parameters
The amplitude A of the sinusoidal wave is shown as the length of the electric
vector at a maximum in the wave.
The period p of an electromagnetic wave is the time in seconds required for the
passage of successive maxima or minima through a fixed point in space is called the
period, p, of the radiation.
In any medium containing matter, propagation of radiation is slowed by the
interaction between the electromagnetic field of the radiation and the bound electrons
in the atoms or molecules present.
Two other important wave parameters are frequency and wavelength.
• The frequency ν of an electromagnetic wave is the number of oscillations that
occur in one second. The unit of frequency is the Hertz (Hz) or S-1.
• The wavelength λ of an electromagnetic wave is the linear distance between
two successive maxima or minima. Its SI unit is metre (m). Other unit of nanomet
(nm, 1nm = 10-9 m) and Angstrong (Ao, 1 Ao = 10-10m) are also used for very short
wavelengths.
• The velocity C of all electromagnetic wave in vacuum is determined to be the
product of its wavelength λ and velocity ν and to be equal to 2.99792x108 m/s (in air
about 0.03% less):
C = λν = 3.00 x 108 m/s
Example: Calculate the frequency of an electromagnetic having the wavelength
of 325 cm.
Solution:
2
ν = C/λ = 3.00 x 108 m.s-1/(325 x 10-2 m) = 9.23 x 107 s-1

 The wavenumber n , which is defined as the reciprocal of the wavelength in
centimeters, is another way of describing electromagnetic radiation. The unit for
wavelength is cm-1.
Wavenumber is widely used in infrared vibrational spectroscopy. Wavenumber is

useful because, in con, it is directly proportional to the frequency (and thus the
energy) of radiation. Thus we can write: n = kν
Where the proportionality constant k is dependent on the medium and equal to
the reciprocal of the velocity of radiation.
Example:
Calculate the wavenumber of a beam of infrared radiation having a wavelength
of 5.00µm.
−
ν = 1/(5.00x10-4) = 2000 cm-1
II.Electromagnetic spectrum
The electromagnetic spectrum of radiation includes a numerous range of
wavelength (and also energy) in which the portion of visible light detected by
naked eyes is very small compared with the other spectrum regions. It is noted that
methods based on electromagnetic spectrum use not only visible light but also
ultraviolet, infrared, and they are called spectrophotometry methods although
naked eyes are insensitive to the latter types.
The following table lists wavelength and frequency of important spectrum
regions for analytical purpose, and also names relevant spectrophotometric
methods. The last column lists types of quantum transition of nucleus, electron,
atom, and molecule that are the foundation of different electromagnetic spectrum
techniques.
Table 0.1: Application of electromagnetic radiation in analysis
Types of Wavelength Wavenumber Types of quantum
electromagnetic range range(cm-1) transition
spectrum techniques
Gamma ray emission 0.005 – 1,4 Ao - Nucleus
Diffraction, 0.1-100 Ao - Nonbonding electrons

fluorescene, emission,
absorption of X rays
10-180 nm 1 x 10 6 đến 5 x 104 Bonding electrons
Vacuum ultraviolet
absorption
180-780 nm 5 x 10 4 đến 1.3 x Bonding electrons
Absortion, emission 104
and fluorescene of
3
ultraviolet and visible

light 0.78-300 µm Harmonic
4
1.3 x10 đến 3,3 x vibration/rotation of
Infrared absorption 0.75-3.75 mm 101 molecule
13-27 Vibration and rotation
Microwave absortion 3 cm of molecule
0.33
Electron spin 0.6-10 m Spin of electron in
-2 3
resonance 1.7x 10 đến 1 x 10 magnetic fields
Spin of nuleus in
Nuclear spin magnetic fields
resonance
Chapter 1: ULTRAVIOLET-VISIBLE (UV-VIS) ABSORPTION

SPECTROSCOPY
1.1. UV-VIS Absorption of molecules

 Molecules under normal conditions exist stably in the lowest energy state
called ground state. When a molecule is provided energy, for instance by a radiation
source having suitable frequency to stimulate, bonding electrons in the molecule
absorb the energy and the molecule is promoted to an exited state.
 According to quantum mechanics, in the ground state of a molecule its
electrons are filled fully into bonding molecular orbitals of σ and π. Apart from these
molecular orbitals, some molecules still have low energy non-bonding valence
electron pairs of n (for instance a nitrogen atom in amine compounds has a non-
bonding valence electron pair).
 If the above molecules are illuminated by a beam of light having suitable
frequency v, low energy existing electrons are stimulated to transfer to higher energy
antibonding molecular orbitals such as σ→σ*, π→π* , n → σ*, and n → π *. This
state is called exited state that is unstable and exists in a very short time interval
(about parts per million of second).
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 Energy of hν of a photon is equal to the difference between the two molecular

orbital energies. The transition of an electron between two orbitals is called electronic
transition, and the absorption process is called electronic absorption.
 In addition to electronic transitions, molecules exhibit two other types of
radiation-induced transition: vibrational transitions and rotational transitions.
Vibrational transitions come about because a molecule has a multitude of quantized
energy levels (or vibrational states) associated with the bonds that hold the molecule
together.
 Figure 1.1 is an energy-level diagram that depicts some of the processes that
occur when a diatomic species of CO absorbs suitable infrared, visible, and ultra-
violet radiation. In this diagram the energy E1, one of the several electronically
excited states of the molecule, is shown relative to the energy of its ground state E o. In
addition, some of vibrational levels associated with each of the electronic states of a
molecule and in turn, some of rotational levels associated with each of vibrational
states are depicted by the lines labeled 1, 2, 3, and 4 in Figure 1.1 (The lowest
vibrational and rotational levels are labeled 0.
 The overall energy E associated with a molecule is then given by:
E = Eelectronic + Evibrational + E rotational
Where: Eelectronic (denoted Ee) is the energy associated with the electrons in the
various outer orbitals of the molecule,
Evibrational (Ev) is the energy of the molecule as a whole due to interatomic
vibrations,
And Erotational (Er) is the energy associated with rotation of the molecule about its
center of gravity.
 Suppose that the energy of a molecule in the ground electronic state and the
excited electronic state is denoted Eo, and En, respectively. So the difference between
these two energy levels is equal to the energy hν of photon, which is absorbed by this
molecule:
ΔE = En – Eo = ΔE(e) + ΔE(v) + ΔE(r) = hν = hc/λ
Here, ΔE(e) > ΔE(v) > ΔE(r)
 As can be seen in Figure 2, the molecular absorption in the ultraviolet and
visible regions consist of absorption bands made up of closely spaced lines, leading to
the fact that a typical UV-visible spectrum consists of a multilude of lines unlike the
spectrum of atomic emission or atomic absoption.
 So the UV-VIS molecular absorption spectrum is derived from the interaction
of valence electrons in a molecule with a light beam having its wavelength in the UV-
VIS region.
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Figure 1.1: Electronic, vibrational and rotational energy levels corresponding to the ground
electronic state and the first excited electron state of diatomic molecules (for example CO 2)
Figure 1.2: UV-VIS molecular absorption spectrum of some organic compounds (continous
spectrum) and atomic absorption spectrum of Na (line spectrum).
1.2. Beer’s law

1.2.1. The law
If we illuminate a beam of parallel monochromatic radiation in UV-VIS regions
with power Io that strikes a cell containing a solution of an absorbing species with a
length l, a significant portion of the beam power is lost by reflection at the two
air/wall interfaces as well as at the two wall/solution interfaces, and the remaining is
absorbed by absorbing species in the solution.
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If after passing throught the cell, its power of the beam is decreased to I 1 as a
result of absorption , the UV-VIS molecular absorption of the species in the cell can
be expressed by the two following terms:
Absorbance A = log(Io/I) or
I
x 100 ( )
Transmittance T = I o
According to Beer’s law: A = εlC

Here:
ε is a proportionality constant called the molar absorbility. The magnitude and
dimention of ε depends on the units used for l, the path length and C, the
concentration of absorbing species. The magnitude of ε also depends on the
wavelength of used monochromic radiation.
For solutions of an absorbing species, b is often given in terms of centimeters,
and C in grams per liter. Absorptivity then has units of L.g-1.cm-1)
1.2.2. The application of Beer’s law to mixtures
Beer’s law also applies to a medium contining more than one kind substance. If
there is no interaction between various species, the total absorbance for a
multicomponent solution is given by the following expression:
A total = A1 + A2 + …. + An = ε1lC1 + ε2lC2 + … + εnlCn
where the subscripts refer to absorbing components 1, 2,…,n .
Consequence: A measured solution = Aanalyte + A blank
Here:
Aanalyte is the absorbance of solution only containing the analyte;
Ablank is the absorbance of blank solution, which is prepared such a way that this
solution have all components in the measured solution except the analyte.
1.2.3. Factors affecting Beer’s law
a. Effect of concentration:
Beer’s law is successful in describing the absorption behavior of media
containing relatively low analyte concentrations.
At high concentrations (usually > 0.01M), the average distance between the
species for the absorption is diminished to the point where each affects the charge
distribution of its neighbors. This interaction, in turn, can alter their ability to absorb
the given wavelength of radiation, leading to causing deviation from the linear
relationship between absorbance and concentration.
A similar effect is sometimes encountered in media containing low absorber
concentrations but high concentrations of other species, particularly electrolytes. The
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close proximity of ions to the absorber alters the molar absorbility of the latter as a
result of electrostatic interactions. The effect is lessened by dilution.
Deviation from Beer’s law also arise because ε is dependant on the refractive
index of the medium. Thus, if concentration changes cause significant alterations in
the refractive index of a solution. However, this effect is rarely significant at
concentrations less than 0.01M.
b. Effects of chemical processes
Apparent deviations from Beer’s law arise when an analyte dissociates,
associates, or reacts with a solvcent to generate a product that has a different
absorption spectrum from that of the analyte.
A typical example is the dilution of a acid-base indicator leading to the shift in
the equilibrium:
HInd  H+ + Ind-
As a result, if the solution is diluted doubly, the concentrations of Hind and Ind-
are not decreased a half respectively. Vì vậy nếu pha loãng gấp đôi thì nồng độ chất
màu Hind không tương ứng giảm đi một nửa. This leads to a deviation from Bee’s
law.
c. Effects of polychromatic radiation

Beer’s law is observed truly only with monochromatic radiation. Howerver, in
practice it is difficult to produce a single wavelength because devices that isolate
portions of the output from a continous source produce a more or less symmetric band
of wavelength around the desired one, λ ± Δλ.
The better monochromator the spectrometer has, the smaller band of wavelength.
This leads to a small deviation from Beer’s law. The deviation can be limited to
insignificant level if the maximum absorption band of wavelength is chosen to
measure sample.
Figure 1.3 illustrates the effect of the absorption leads to a large deviation
polychromatic radiation on Beer’s law. from Beer’s law due to significant
If band 1 is chosen to measure sample, change of ε in this band, so the
the deviation from Beer’s law is relationship between A and C is no
insignificant (Figure 1.3b shows the longer linear.
linear relationship between absorbance
and concentration) because the molar
absorbility ε does not change
significantly over the whole band.
Meanwhile, using band 2 to measure
8
Figure 1.3: The effect of polychromatic

radiation on Beer’s law.
d. Effect of the stray radiation
The radiation used for absorbance measurements is usually contaminated

with small amunt of stray radiation due to instrumental imperfections. Stray
radiation is the result of scattering phenomena off the surfaces of prisms,
lenses, filters and windows.
It often differs greatly in wavelength from the principal radiation and, in addition,
may not haved passed through the sample or solvent, but it may enter
photomultiplier leading to a deviation from Beer’s law.
1.3. Measuring principle and instrument components for UV-VIS spectrometry
1.3.1. Measuring principle
UV-VIS spectrum is molecular absorption one of an analyte in the homogeneous
state with one of such solvents as water, methanol, benzene, acetone, CCl 4 and so on.
So measuring steps may be as follows:
Firstly, If the analyte has a sensitive absorption spectrum in ultraviolet-visible
regions, it can be dissolved in a suitable solvent. Otherwise, a reaction of the analyte
(for example metal ions) with an excessConcentration
reagent in a suitable solvent is performed to
form a compound having a sensitive UV-VIS spectrum.
The next is to illuminate the solution containing the analyte or its product by a
monochromatic radiation beam having suitable wavelength so that the analyte or its
product absorbs the radiation to produce its UV-VIS spectrum. Therefore, the solution
must be contained in a cell having a determinate length.
The last is to obtain the spectrum, chose a desired wavelength at which the
absorbance A is maximum, meaning that the power of the beam is measured before
and after passing through the solution.
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Based on such measuring principle, a UV-VIS spectrometer needs to have the

following fundamental components.
1.3.2. Fundamental instrument components
a. Cuvettes (cells)
A cuvette is a small tube of square cross section, sealed at one end, made of
plastic, glass, or fused quartz (for UV light) and designed to hold samples for UV-VIS
absorption measurements. There are three types of cuvettes: plastic, glass and quartz
cuvettes. Disposable plastic cuvettes are often used in fast spectroscopic assays, where
speed is more important than high accuracy. While plastic and glass cuvettes are only
for use in the wavelength range of visible light, fused quartz cuvettes can be used in
the UV through visible ranges. It is noted that plastic cuvettes are not sutable for
measuring samples with organic solvents.
Figure 1.4: Some types of UV-VIS cuvettes

b. Radiation sources
 Radiation sources for molecular absorption must generate a beam with
sufficient power in the wavelength regions to permit readily detection and
measurement. In addition, the source must be stable during the time of measuring so
that the measuring results can be reproducible.
 The most common source of visible and near-infrared is the tungsten-filament
lamp generating the wavelength region between 380 nm đến 2500 nm.
 Tunsten – halogen lamps contain a small quantity of iodine within a quartz
envelope that houses the filament, which can operate at a temperature of about 3500
K. This leads to higher intensities and extends the range of the lamp well into the
ultraviolet. This lamp also has the longer lifetime because the sublimation of the
tungsten from the hot filament is limited. In the presence of iodine, sublimed tungsten
reacts with gasous iodine to form gasous WI 2 molecules, which then diffuse back to
the hot filament, where they decompose and redeposit as tungsten atoms.
 Deuterium D2 lamp
Most modern lamps for generating ultraviolet radiation contain deuterium
and are of a low voltage type in which an arc is formed between a heated,
oxide-coated filament and a metal electrode. The heated provide electron to
maintain a direct current at a potential of about 40 V; a regulated power supply
10
is required for constan intensities. This lamp provides a useful continuous

spectrum in the region from 160 nm to 375 nm.
c. Wavelength selectors
 Radiation filters: Filters operate by absorbing all but a restricted band of
radiation from a continuous source. There are two types of filters used in absorption
measurements: interference filters and absorption filters. While interference filters can
use with both ultraviolet and visible radiation, absorption filters are limited in
application to the visible region.
 Monochromators: As can be seen in Figure 1.6, there are two general types of
monochromators, one of which uses a grating to disperse radiation into its component
wavelength; the other uses a prism for this purpose.
 One of the most common types of reflection gratings is the Echellete grating,
which is depiced in Figure 1.5. The length of the grating is from 3 to 10 cm, and its
hard, optically flat, posished surface is ordinalily made by a coating of aluminum, or
sometimes gold or platinum. On the surface a large number of parallel and closely
spaced grooves have been ruled with a suitable shaped diamond tool. A grating for the
violet and visible region typically contains from 300 to 2000 grooves/mm, with 1200
to 1400 being most common.
 In Figure 1.5, a parallel beam of monochromatic radiation is shown
approaching at an angle I relative to the grating normal. The incident radiation
consists of three parallel beams labeled 1,2,3. The diffracted beam is reflected at the
angle r, which depends on the wavelength of the radiation.
Figure 1.5: Mechanism of diffraction from an echellette-type grating
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Figure 1.6: (a) Two types of chromators (a) grating monochromator (b) prism monochromator
d. Photo Detectors
 A photo detector is a device that converts light intensity into electric signals
that can be subsequently ampflified, manipulated, and finally converted into numbers
that are related to the magnitude of the original signal.
 There are many types of photo detectors in which two most common types
used for UV-VIS spectrometers are phototubes and photomultifliers.
Phototubes:
 As can be seen in Figure 1.7a, a phototube consisits of a semicylindrical
cathode and a wire anode sealed inside a evacuated transparent envelope. The concave
surface of the cathode supports a layer of photosemissive material such as a alkali
metal or metal oxide that tends to emit electrons upon being irradiated. When a
potential is applied across the electrodes, the emitted photoelectrons flow to the wire
anode, producing a current that is readily amplified and displayed or recorded.
 The number of electrons ejected from a photoemissive surface is directly
proportional to the radiant power of the beam striking that surface.
Photomultiflier tubes (PMT):
 As shown in Figure 7b, the photomultiflier tube (PMT) is similar in
construction to the phototube but is significantly more sensitive. Electrons emitted
from the cathode are accelerated toward a dynode (labeled 1 in the figure) maintained
at a potential 90 V more positive than the cathode. When striking the dynode surface,
each accelerated photoelectron produces several additional electrons, all of which are
then accerelated to dynode 2, which is 90 V more positive than dynode 1. Here again,
electron amplitication occurs. By the time this process has been repeated at each of the
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remaining dynodes, 106 to 107 electrons have been produced for each photon; these
electrons are finally collected at the anode. So the resulting current is then further
amplified electronically and measured.
(b)
Figure 1.7: (a) Diagram of a phototube and (b) a photomultiflier tube (PMT)
e. Designs of UV-VIS spectrometers
Single beam instrument:
Figure 1.8a shows the design of a typical single beam spectrometer, and Figure
6b shows its optical diagram. It is simple, easy to maitain and inexpensive, so it is a
very common in use throughout the wolrd.
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Figure 1.8: a) a typical single beam spectrometer, b) its optical diagram

Double –beam instruments:
Figure 1.8bis: Structural diagram of a double beam spectrometer.
1.4. Quantitative analysis applications

 UV-VIS molecular absorption spectroscopy is one of the most useful tools for
quantitative analysis. The most important characteristics of spectrometric and
photometric methods are:
 Wide applicability: Absorption spectroscopy can be used to determine
quantitatively and directly enormous numbers of inorganic, organic, and biochemical
compounds that absorb ultraviolet or visible radiation. In the case of many
nonabsorbing species, they can also be determined after conversing them to absorbing
derivatives by chemical reactions.
 High sensitivity: Typical detection limits for absorption spectroscopy range
from 10-4 to 10-5M. This range can be extended to 10 -6 or even 10-7 M with certain
procedural modification.
 Moderate to high selectivity: If a wavelength can be found at which the analyte
alone absorbs, preliminary separations become unnecessary. In the case of occurring
overlapping absorption bands, corrections based on additional measurements at other
wavelength sometimes eliminate the need for a separation step.
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 Good accuracy: The relative error in concentration for a analytical procedure

based on UV –VIS absorption spectroscopy lie in the range from 1% to 5%. The
measurements are easily and rapidly perdormed with modern instruments.
1.4.1. Calibration curve method (External standard method)
In most spectrophotometric methods, calibration is achieved by the
method of external standards. Here, a series of standard solutions of the analyte
is used to construct a calibration curve of absorbance versus concentration or to
produce a linear regression equation.
 The slope of the calibration curve or regression equation is the product of
absorptivity and pathlength. Thus, using external standards is way of determining the
proportionality factor between absorbance and concentration under the same
conditions and with the same instrument as is used for the samples.
1.4.2. Standard Addition Method
1.4.2.1. Principle
 The difficulties associated with production of standards with an overall
composition closely resembling that of the sample can be formidable. Under such
circumstances, the method of standard additions may prove useful.
 In the single-point standard addition method, a known amount of analyte is
introduced into a second aliquot of the sample and the difference in absorbance is used
to calculate the analyte concentration of the sample.
 Alternatively, multiple additions can be made to several aliquots of the sample
and multiple standard addition calibration curve obtained.
1.4.2.2. Procedure for using many increments of a standard solution
 Several identical aliquots Vx of the unknown solution with a concentration Cx
are transferred to volumetric flasks having a volume Vt. To each of these flashs is
added a variable volume Vs mL of a standard solution of the analyte having a known
concentration Cs. The color development reagents are then added, and each solution is
diluted to volume Vt.
 If Beer’s law is followed, the absorbance As of the solutions is described by
As = εlVsCs/Vt + εlVxCx/Vt = kVsCs + kVxCx
= mVs + b
where k is a constant equal to εl/Vt ; m =kCs and b = kVxCx
 A least-squares analysis can be used to determine m and b, and then C x can be
obtained rong mẫu như sau:
m/b = kCs/kVxCx
Cx = bCs/mVx
Example:
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 Use a suitable pipette to add 10.0 mL of a water sample to to each of 5

volumetric flasks of 50.00 mL, and then add to each of the flasks a precisive 11.1
ppm Fe3+ standard volume of 0.00; 5.00; 10.00; 15.00 and 20.00 mL, respectively.
 Add thiocyanat SCN- to form the red complex and dilute to the volume of each
of the flask. The absorbances of the solution of each flask are measured as follows:
0.240; 0.437; 0.621; 0.809 and 1.009. Determine the concentration of Fe 3+in the
sample.
Ans: A = 0.03820 Vs + 0.2412  CFe3+ = 7.01ppm
1.4.2.3. Procedure for using two increments of a standard solution
In the interest of saving time or sample, it is possible to perform a
standard addition analysis using only two increments of sample.
Here, a single addition of Vs mL of standard would be added to one of the
two samples and we can write
A = εlV C /V
1 x x t
A = εlV C /V + εlV C /V
2 s s t x x t
Where A1 and A2 are absorpbances of the diluted sample and the diluted
sample plus standard, respectively. Dividing the second equation by the first
gives
C = A C V /(A -A )V
x 1 s s 2 1 x
Examle: A 2.00 mL urine the milligrams of phosphate per
specimen was treated with milliliter of the specimen
reagent to generate a colored (mg/ml). Ans : 0,292
adduct with phosphate, and the
sample was then diluted to 100
mL. %.00 mL of a phosphate
solution containing 0.0300 mg
phosphate/mL were added to a
second 2.00 mL sample, which
are treated in the same way as
the orogonal sample. The
absorbance of the first solution
was 0.428 and 0.538. Calculate
1.4.3. Analysis of Mixtures

The total absorbance of a solution at any given wavelength is equal to the
sum of the absorbances of the individual components in the solution. This
relationship makes it possible in principle to determine the concentration of the
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individual components in a mixture even if there is strong overlap in their

spectra. There is no wavelength at which the absorbance is due to just one of
these components. To analyze the mixture, molar absorptivities are first
determined at wavelengths λ1 and λ2. The wavelengths selected are ones at
which the two spectra differ significantly. Thus, at λ 1, the molar absorptivity of
component M is much larger than that for component N. The reverse is true for
λ2. To complete the analysis, the absorbance of the mixture is determined at the
same two wavelengths. From the known molar absorptivities and pathlength,
the following equations hold:
A1 = εM1bcM + εN1bcN
A2 = εM2bcM + εN2bcN
.
Problems
1.1) Calculate the energy of a mole of photons corresponding to a wavelength of
300 nm.
1-2)
1-3) Calculate the absorbance of an organic dye (C =7×10 −4mol L−1), knowing hat the
molar absorptivity = 650mol L−1cm−1and that the length of the optical path of
the cell used is 2×10−2m. What would happen to the absorbance if the cell used
was of double its present thickness?
1.4) A solution containing 6.23 ppm KMnO4, had a transmittance o[ 0.195 in a 1.00-
cm cell at 520 nm. Calculate the molar ahsorptivity of KMnO4, at 520 nm.
1.5) A solution containing 5.24 mg/100 mL of A (335 g/mol) has a transmittance of
55.2% in a 1.50 cm cell at 425 nm. Calculate the molar absorptivity of A at this
wavelength.
1.6) A 5.00 × 10–4 M solution of an analyte is placed in a sample cell that has a
pathlength of 1.00 cm. When measured at a wavelength of 490 nm, the
absorbance of the solution is found to be 0.338. What is the analyte’s molar
absorptivity at this wavelength?
1.7) A compound X is to be determined by UV-visible spectrophotometry. A
calibration curve is constructed from standard solutions of X with the
following. results: 0.50 ppm, A = 0.24; 1.5 ppm, A = 0.36; 2.5 ppm, A = 0.44;
3.5 ppm, A = 0.59; 4.5 ppm, A = 0.70. A solution of unknown X concentration
had an absorbance of A = 0.50. Find the slope and intercept of the calibration
17
curve, the concentration of the solution of unknown X concentration. Construct

a plot of the calibration curve and determine the unknown concentration by
hand from the plot. Compare it to that obtained from the regression line.
1.8) A 1.28×10−4 M solution of potassium permanganate has a transmittance of 0.5
when measured in a 1 cm cell at 525 nm.
1. Calculate the molar absorptivity coefficient for the permanganate at this
wavelength.
2. If the concentration is doubled what would be the absorbance and the
percentage transmittance of the new solution?
1.9) To determine the concentrations (mol/L) of Co(NO3)2 (A) and Cr(NO3)3 (B) in an
unknown sample, the following representative absorbance data were obtained.
Measurements were made in 1.0 cm glass cells.

1. Calculate the four molar absorptivities: εA(510), εA(575), εB(510) and εB(575).
2. Calculate the molarities of the two salts A and B in the unknown.
1.10) Jones and Thatcher developed a spectrophotometric method for analyzing
analgesic tablets containing aspirin, phenacetin, and caffeine. The sample is
dissolved in CHCl3 and extracted with an aqueous solution of NaHCO 3 to
remove the aspirin. After the extraction is complete, the chloroform is then
transferred to a 250-mL volumetric flask and diluted to volume with CHCl 3. A
2.00-mL portion of this solution is diluted to volume in a 200-mL volumetric
flask with CHCl3.
The absorbance of the final solution is measured at wavelengths of 250 nm and
275 nm, at which the absorptivities, in ppm–1 cm–1, for caffeine and
phenacetin are
caffeine: ε250 = 0.0131 and ε 275 = 0.0485
phenacetin: ε 250 = 0.0702 and ε 275 = 0.0159
Aspirin is determined by neutralizing the NaHCO 3 in the aqueous solution and
extracting the aspirin into CHCl3. The combined extracts are diluted to 500 mL
in a volumetric flask. A 20.00-mL portion of the solution is placed in a 100-mL
volumetric flask and diluted to volume with CHCl 3. The absorbance of this
solution is measured at 277 nm, where the absorptivity of aspirin is 0.00682
ppm–1 cm–1. An analgesic tablet treated by this procedure is found to have
absorbances of 0.466 at 250 nm, 0.164 at 275 nm, and 0.600 at 277 nm when
using a cell with a 1.00-cm pathlength. Report the milligrams of aspirin,
caffeine, and phenacetin in the analgesic
18
tablet.
Chapter 2: ATOMIC EMISSION SPECTROSCOPY (AES)

2.1. The appearance of AES
In order to understand the appearance of AES, we look firstly into the
structure of an atom. This is due to the fact that AES is the spectroscopy of
excited gaseous unbound atoms, and the structure of the atom is a key factor of
the appearance of AES. The atomic structure consists of two main parts:
- The atomic nucleus: it decides positively charged part and the mass of
an atom. The nucleus lies at the centre of the atom, but it occupies a very small
volume (1/1000) compared to the whole volume of the atom (Figure 2.1).
Figure 2.1: A simple scheme of electron layers and their energies

- Electron shells: Electrons in each a round the nucleus, the smaller
shell of an atom orbit in different energy they have, and according to
orbitals around the atomic nucleus. the Klechkowski rule, their energies
The more closely the electrons orbit increase gradually in the order of
19
the arrangement of the shells from

inside to outside as can be seen in
Figure 2.2. Electrons occupy a very
large volume of the atom. The
valence electrons in the outermost
shell of atom are capaple of forming
both emission and absorption
spectra.
Figure 2.2: Energy scheme of subshells
Under normal conditions, atoms existing in matter are at the stable state, which
has the minimum energy, called fundermental state of matter. These atoms neither
absorb nor emit energy. However, if the matter is evaporated to transfer the atoms to
free atomic gaseous state, and these atoms are provided suitable energy, they are
excited such that their outermost electrons will be transferred to other orbitals having
higher energy:
Ao + ΔE → A*
(Ground) (Excited)
Here, Ao and A* are ground and excited electronic states of the atom respectively;
DE is the provided energy that is smaller than the first ionization energy of the atom.
These excited atoms are in 2.3). So an emission spectrum is one

unstable state and only exist at this that is formed by transferring energy
state in a very short time (about 10 -8 s), levels of the atom.
and then they tend to release the A* ��
� A 0 + n(hg )
received energy until they return to the
A beam of the emission
stable initial ground state. That is the
radiation has n monochromatic rays
emission process of the excited
with different wavelengths (mainly in
electronic atoms. The emission
the UV-VIS region), and here, n is an
radiations of this process form the
integer number from 1 to thousands.
emission spectrum of atom (Figure
20
The more outmost electrons the

atom of an element has, the greater n
value, meaning that the more emission
lines.
Figure 2.3: Emission spectrum of Na
According to the Einstein equation, we have:

hg = DE
or:
h.c
DE =
l
Here,
c is the light velocity in vacuum;
h is the plank constant;
l is the wavelength of the emission line.
If the emission beam is collected and dispersed into its component
wavelengths, an emission spectrum with the whole emission lines of the excited
element is obtained. So AES spectrum is a product of interacting between matter (free
atoms) and a suitable energy source such as electricity or thermal energy, which is
used to stimulate the atoms to form emission spectrum. Now, if the emission intensity
of a monochromatic light (an emission line) is called I , under specific conditions, we
have the following expression:
Il = k.Cb (2.1)
Here,
k is a experimental constant dependent on such many factors as conditions
used to evaporize sample and to promote some atoms to excited electronic atoms. If
these conditions are fixed, k can be constant.
21
b is also constant, and this value is dependant on the nature of each atom, each
spectrum line, and the concentration of each analyzing element. The value of b always
lies in the following region:
0<b≤1
Here, for every spectrum lines, we always have a specific value of the
concentration C0 of analyte in sample at which the value of b is equal to 1, and we
always have:
- If the concentration of the analyst is smaller than the value of C 0, the value of
b is always equal to 1.
- Otherwise, if the above concentration is higher than C 0, the value of b
gradually decreases from 1 to zero with the increase in the concentration of the
analyte.
Formula (2.1) is a fundermental equation of the quantitative analysis method
based on atomic emission spectroscopy, and the relationship between emission
intensity Il and concentration C is illustrated by a plot in Figure 2.4.
Iλ
C
B
0 Ca Co C
Figure 2.4: Relationship between the intensity Iλ of spectrum line and concentration C
As can be seen in the plot, the curve has two sections. Section AB corresponding
to the concentration range of Ca-C0 is linear. Here, the relationship between Il and C
is linear (like equation y=ax). Otherwise, section BC is not linear. The concentration
C0 is called the limit concentration of linear region. The more sensitive the spectrum
22
line, the greater the value of C0 and vice versa. In analytical practical, section AB is
always used due to easier construction of calibration curve and results derived from
the curve having higher precision.
Apart from quatitative analysis, atomic emission spectroscopy is a super
qualitative analysis method in detecting the presence of many elements, especially
metals, in different samples. That is due to the fact that each element has some distinct
spectrum lines that other elements do not have.
2.2. Effect of temperature on atomic spectrum
2.2.1. The Boltzmann Distribution
Consider an atom with energy Eo at ground state. If this atom absorb energy from
the medium, it transfer to excited state with energy E*. In general, an atom (or
molecule) may have more than one state at a given energy. The number of states at
each energy is called the degeneracy. We will call the degeneracies go and g*. The
Boltzmann distribution describes the relative populations of different states at thermal
equilibrium. At a given temperature, if equilibrium exists, the relative population
(N*/No) of any two states is
N ∗¿ g∗¿ − ΔE/ kT
= e
No go
¿
¿
where T is temperature (K) and k is Boltzmann’s constant (1.381x10-23 J/K).

2.2.2. Effect of temperature on excited-state population
The lowest excited state of a sodium atom lies 3.371x10 -19 J/atom above the
ground state. The degeneracy of the excited state is 2, whereas that of the ground
state is 1. The fraction of Na in the excited state in an acetylene-air flame at 2 600
K is, from the above equation,
N∗¿ 2 −( 3 .371×10−19 J / 1. 381×10−23 J / K )×(2600 K )
= e =1 .67×10−4
No 1
¿
That is, less than 0.02% of the atoms are in the excited state. If the temperature
were 2 610 K, the fraction of atoms in the excited state would be
N * 2 - (3.371�10-19 J /1.381�10-23 J / K )�(2610 K )
= e = 1.74 �10 -4
No 1
The fraction of atoms in the excited state is still less than 0.02%, but that fraction
has increased by 100(1.74 - 1.67) /1.67 = 4%.
2.2.3. Effect of temperature on emission and absorption
As can be seen from the above calculation, at 2600 K there are more than
99.98% of sodium atoms are in their ground state. Varying the temperature by 10 K
23
hardly affects the ground-state population and would not noticeably affect the signal
in atomic absorption. Emission intensity is proportional to the population of the
excited state. Because the excited-state population changes by 4% when the
temperature rises 10 K, emission intensity rises by 4%. Almost all atomic emission is
carried out with an inductively coupled plasma, whose temperature is more stable than
that of a flame.
2.3. Measurement principle and main components of an atomic emission
spectrometer
2.3.1. Measurement principle
Based on the appearance of atomic emission spectrum as mentioned above, for
applying the emission spectroscopy to analyze elements, we need to perform the
following steps:
1. Find suitable conditions to vaporize completely the sample.
2. The gaseous-state sample is completely atomized to form ground–state
atoms of the analyzing element.
3. Excite these atoms to generate emission spectrum in such a way that the
exciting efficiency is high, stable and reproducible.
4. Collect the whole emission beam of the sample, disperse this beam to obtain
emission spectrum of the sample/
5. Based on this spectrum to analyze qualitatively and quantitatively.
These five steps have to be performed for an analytical procedure using AES to
determine an element in a specific sample in which each step is standardized, and the
most suitable conditions are chosen for each step.
II.3.2. Main components
According to the general principle, an atomic emission spectrometer has to
consist of three main components as follows:
1. Power source: it is a divice or a component to perform steps of 1-3,
vaporization, atomization and exciting to form emission spectrum.
2. Spectrophotometer: it is a component containing monochromator and
photomultipliertube that has duties of collecting, dispersing the emission beam into
discrete spectrum lines of each element and measuring the intensities of these
spectrum lines.
3. Readout device: it displays the results of qualitative and quantitative
analysis.
These above components are indispensable for an atomic emission
spectrometer. In the present, in order to programme measuring processes, increase
24
measuring rate, or store, print, compare results, an atomic emission spectrometer is

coupled some more components such as autosampler, computer.
2.4. Energy sources
This component is very important because it decides the sensitive of
measurement due to the fact the sample needs energy to convert into vapour state of
atoms, and then excited to form emission spectrum of analyzing element. As a result,
the energy source needs to be sastified the following conditions:
- The source has energy enough great to evaporize the sample, atomize them,
and excite these atoms with high efficiency.
- It has to be stable and reproducible,
- It can adjust the magnitude of energy that is suitable for measuring each
element in a specific sample.
- It does not generate its spectrum that affects or makes difficult for the
analysis.
- It should comsume the sample not so much.
Among these requirements, the fifth is not necessary in some cases such as soil
analysis, but it is important when analyzing blood, serum, or super pure analysts.
Due to the satisfaction of these above requirements, although there are a lot of
energy sources, only some types such as flame, arc, laser, plasma have been used in
AES. Afterward, let us consider one by one their features that are commonly using in
AES.
2.4.1. Flames
The flame is formed when igniting a mixture of an oxidant (such as O 2, air, or
N2O) and a fuel (such as C2H2, propane) to generate the temperature of the flame. The
flame temperature is the key factor for the process of evaporating, atomizing and
exciting to produce the emission spectrum of the sample. In general, the flame could
be generated the temperature of from 1700 to 33000C.
Table 2.1. The flame temperature of gaseous mixtures.
Oxidant Fuel Temperature

Ratio
(K)
Air Acetylene 4,2/0,9 2100
N 2O Acetylene 4,0/4,5 2800
Air Propane 4,2/1,5 2050
Oxygen Acetylene 1,2/1,5 2700
25
The flame temperature depends upon:

- the structure of premix burner in which fuel, oxidant, and sample are mixed
before introduce into the flame,
- the nature of oxidant and fuel used to generate the flame,
- the ratio of oxidant and fuel.
Consequently, a mixture of two gases (one oxidant and one fuel) could also
generate flames having different temperatures when changing the rato of these two
gases.
For flames having low temperature, they are only suitable to exciting atoms of
alkali metal or alkaline earth metal because these metals have low emission
exciting energy (2 – 3.5 eV).
Flame consists of three regions as

be shown in Figure 2.5. These are the
primary combustion zone, interconal
region and outer cone (secondary
combustion zone). The appearance and
the relative sizes of these regions can
be changed with the fuel-oxidant ratio.
The primary combustion zone of the
flame is blue in color. In this region,
there is no thermal equilibrium.
Therefore, it is not used in flame
spectroscopy. The interconal region is
rich in free atoms and is the most
widely used region for the
spectroscopy. In The outer cone the
Figure 2.5: The structure of a flame products of the inner core are converted
to stable molecular oxides.
Due to low flame temperature, flames have been only used as an exciting energy
source to analyse alkali metals (Li, Na, K, Rb, Cs) and some alkaline earth metals (Ca,
Mg) with the sensitivity of 1 - 10 µg/ml. These metals are determined using a method
called Flame Spectrophotometry, but in reality it is based on the principle of atomic
emission spectroscopy. Here, it should be noted that the spectrum excitation by flames
is subjected to the considerable effect of sample matrix.
26
Figure 2.6: Diagram of a flame spectrometer using to determine alkali metals

2.4.2. Electric Arc
The usual arc source for a spectrochemical analysis is formed with a pair of
graphite or metal electrodes spaced a few millimeter apart. This is a thermal energy
source generated by the discharge between two electrodes (mainly grafite) applied
high current (10 - 20A) and average potental (220 - 250V). The temperature of arcs is
about 3000 - 60000C, so arc is a moderate power source. Its temperature depends on:
- materials used to be made of the electrode,
- electric current flowing though the arc circuit,
- and the matrix composition of sample.
In general, if materials used to make the electrode is thermally stable, having
high ionization and excitation potentials, the arc would have high temperature, and the
larger electric current, the higher arc temperarure.
As a result of these above features, arc is a suitable excitation source to generate
spectra for 45 elements. It is also an energy source of forming an analysis having the
relatively high sensitivity (about 10 – 0.1 µg) of atomic emission spectroscopy.
Electric arc has two types of alter current arc and direct current arc in which the
former gives higher stability leading to being used more than the latter. This type can
be used to analyze either electrically conductive or resitive power samples. It should
be noted that atomic emission spectrum using electric arc is seriously affected by the
matrix of sample, especially the thermally stable ones such as silicate, zircomate,
wolfamate.
Table 2.2. Temperature and materials used to make electrodes
Temperature Ionization
Electrodes Current (A)
(0C) potential (eV)
1. pressed 10 5800 (1) 11.25
27
Graphite
2. Fe metal 10 4400 (2) 7.86
3. Zn metal 10 5200 (3) 9.40
4. Al metal 10 3800 (4) 5.90
5. Cr metal 10 4000 (5) 6.76
2.4.3. Spark
This is also a thermal energy source produced by the electric discharge between
two electrodes having small electrical current (100 - 1000 mA) but very high potential
(20 - 30 kV). The discharge is periodically performed with from 50 to 500 times per
second. The electric spark temperature is about 4000 - 7000 0C. Because the spark
does not make the electrodes so hot, the evaporation of sample is difficult and slow.
Consequently, the analytical sensitivity is not so high (10 - 100 µg), but the stability is
better than arc. Therefore using spark leads to long excitation time. The spark
temperature depends mainly upon current density and is affected by the structure of
material used to make the electrodes.
The higher thermal stability of the material, the less effect on the temperatue of
spark. Because the electrodes are not subjected to be red-hot, spark excitation is very
suitable to analyze samples of metal, alloy and solution rather than ore, soil, salt,
oxide, silicate.
2.4.4. Inductively coupled plasma (ICP) source
Plasma is a conducting gaseous
misture containing a significant
concentration of cations and electrons.
In the argon plasma employed for
emission analyses, argon ion and
electrons are the principle conducting
species although cations from the
sample also contribute. Argon ion, once
formed in a plasma, are capable of
absorbing sufficient power from an
external source to maintain the
temperature at up to 10000K.
Figure 2.7: Inductively coupled plasma source
28
In reality ICP is a thermoelectric energy source of a flame maintained by an

induction coil powered by a radio-frequency generator. Figure 2.7 is a schematic
drawing of an inductively coupled plasma torch.
It consists of three concentric quartz tubes through which streams of argon flow
at a total rate of between 11 and 17 L/min. the diameter of the largest tube is about 2.5
cm. Surrouding the top of this tube is water-cooled induction coil a radio-frequency
generator capable of producing 2 kW of energy at about 27 MHz. Inonization of the
flowing argon is initiated by a spark from a Tesla coil. The resulting ions and their
associated electrons, then interact with the fluctuating magnetic (labeled H in Figure
2.7) produced by the induction coil I. This interaction causes the ions and electrons
within the coil to flow in the closed annular paths depicted in the figure. Accerelated
electrons collide with other atoms and transfer their energy to these atoms. Once this
process initiates, electrons will absorb energy from the testla coil to maintain the
temperature from 6000 K to 10000 K in this plasma. Due to very high temperature,
the plasma torch is overheating, so it should be protected by a tangential argon flow to
cool the inside walls of the central tube and to center the plasma radially
Figure 2.8: (a) Photograph and (b) temperature profile of a typical inductively coupled plasma
The aerosol of a liquid sample is generated by two ways: using a conventional
nebulizer with the carrier gas of argon and using ultrasonic nebulizer (being used
nowadays).
29
In the method of using ultrasonic nebulizer, sample solution is directed onto a

piezoelectric crystal oscillating at 1 MHz. The vibrating crystal creates a fine aerosol
that is carried by a stream of Ar through a heated tube, where solvent evaporates. The
stream then passes through a refrigerated zone in which solvent condenses and is
removed. Analyte reaches the plasma flame as an aerosol of dry, solid particles.
Plasma energy is not needed to evaporate solvent, so more energy is available for
atomization. Also, a larger fraction of the sample reaches the plasma than with a
conventional nebulizer.
Figure 2.9: A conventional nebulizer with carrier argon
Figure 2.10: An ultrasonic nebulizer
30
ICP has been using popularly because of the following good points:
- It has very high analytical sensitivity, about 10 – 0.1 ng;
- This method is capable of analyzing thermally stable elements as well as
sample having thermally stable matrix;
- The results have small error and good reproducibility;
- Linear range is wide;
- The rate of analyzing sample is fast (50 - 120 samples/h);
- The effect of sample matrix is insignificant.
Due to these above excellences, ICP excitation is mainly used in atomic
emission spectroscopy compared with other source excitations. Especially, ICP is a
very suitable source excitation to analyze rare-earth metals.
2.5. Monochromator and detector in ICP-AES
In Figure 2.11, the most common form of a monochromator (a Rowland circle)
and detector (photomultiplier; PMT) is shown:
The Rowland system utilizes a concave Echellette-style grating monochromator
to separate the various emission lines and simultaneously focus individual
wavelengths on to a series of slits, with each slit aligned to allow a specific
wavelength of radiation to pass to a detector. The standard detector, a photomultiplier
tube (PMT), was discussed in Section 1.3.2. Some systems use multiple PMTs at
fixed locations to monitor each wavelength simultaneously (Figure 2.11) whereas
other systems use a single PMT and move it to different locations to detect each
wavelength. Data from these detectors are processed by a computer because multiple
wavelengths are measured in an ICP-AES system at the same time.
31
Figure 2.11: Overview of a Basic Inductively Coupled Plasma—Atomic Emission

Spectrometry (ICP-AES) from the 1990s.
2.6. Data Collection

When single element detection was the goal of instrumentation, a simple chart
recorder or digital display was sufficient for data collection. However, with the
advances in ICP-AES technology today, computers are a necessity, not only to control
the automatic sampler and the instrument control but also for data collection. Data
collection systems are divided into a “method file” to run the instrument and
“sequence file” to tell the instrument where a sample is in the sample tray, when to run
it, and where to store the collected data file.
2.7. Quantitative analysis by AES
Usually, the determination is done using one of the following methods. In the
determination, the interference and background should be corrected.
2.7.1. Calibration Curve Method
Prepare standard solutions of three or more different concentrations, measure
the emission intensities of these standard solutions, and prepare a calibration curve
from the obtained values. Then, measure the emission intensity for the test solution
with a concentration adjusted to a measurable range, and determine the amount
(concentration) of the object element from the calibration curve.
32
2.7.2. Standard Addition Method

To equal volumes of three or more test solutions, add to each the standard
solution so that the stepwise increasing amounts of the object element are contained in
the solutions, and add the solvent to make a definite volume. Measure the emission
intensity for each solution, and plot the amounts (concentrations) of added standard
object element on the abscissa and the emission intensities on the ordinate on th
e graph paper. Extend the calibration curve obtained by linking the plots, and
determine the amount (concentration) of object element from the distance between the
origin and the intersecting point of the calibration curve on the abscissa. This method
is applicable only when the calibration curve drawn as directed in section 2.7.1 above
is a straight line passing through the origin.
2.7.3. Internal Standard Method
Prepare several solutions containing a constant amount of the specified internal
standard element, and known graded amounts of the standard object element. For
these solutions, measure the emission intensities of the standard object element and
internal standard element at the analytical wavelength of each element under the same
measuring conditions, and obtain the ratios of each emission intensity of standard
object element to the emission intensity of the internal standard element. Prepare a
calibration curve by plotting the amounts (concentrations) of standard object element
on the abscissa and the ratios of emission intensity on the ordinate. Then, prepare the
test solutions, adding the same amount of internal standard element as in the standard
solution. Proceed under the same conditions as for preparing the calibration curve,
obtain the ratio of the emission intensity of standard object element to that of internal
standard element, and determine the amount (concentration) of the object element
from the calibration curve.
2.8. Applications of AES
The large number of elements measurable by AES renders this method of
analysis indispensable since it permits the measurement of elements chosen in
advance or the identification of elements present in any sample. Beyond industrial
analytical applications – for example, monitoring the wear of car and jet engine
motors, without disassembling them, by measuring the metals present in the oil – it is
in the environmental area, without doubt, that analyses are the most widespread.
Applications to the analyses of vegetal or animal production (meat, milk), water,
air (dust and ash discharges by incinerators), or of soils in which elements are present
over a wide range of concentrations (spreading of industrial sludges on agricultural
land) are also common. This methodology has equally a variety of applications in the
forensic sciences or in clinical medicine (analysis of tissues or biological fluids). The
33
advantage of AES is the linearity of response over a very large range of

concentrations, which enables the treatment of complex matrices with a minimum of
preparation. It is frequent that a single solution allows to pass from the measurement
of one highly concentrated element to another, present only as a trace.
Problems
2-1) Li was determined by atomic emission with the method of standard addition.
Use a graph to find the concentration of Li and its uncertainty in pure
unknown. The Li standard contained 1.62 µg Li/mL.
2-2) To measure Ca in breakfast cereal, 0.5216 g of crushed Cheerios was ashed

in a crucible at 600°C in air for 2 h. The residue was dissolved in 6 M HCl,
quantitatively transferred to a volumetric flask, and diluted to 100.0 mL.
Then 5.00-mL aliquots were transferred to 50-mL volumetric flasks. Each
was treated with standard Ca 2+ (containing 20.0 µg/mL ), diluted to volume
with H2O, and analyzed by flame atomic absorption. Find wt% Ca in
Cheerios.
2-3) The potassium concentration in a blood serum is to be analysed using a
method of addition and flame emission. Two extractions of 0.5mL of serum
were taken to create two identical solutions and then both were diluted
further with distilled water to a final volume of 5mL. 10µL of 0.2M KCl is
introduced to one of these. The values obtained from the apparatus were 32.1
and 58.6 arbitrary units. What is the potassium concentration of the serum?
2-4) The concentration of Na in plant materials may be determined by flame
atomic emission. The material to be analyzed is prepared by grinding,
homogenizing, and drying at 103°C. A sample of approximately 4 g is
transferred to a quartz crucible and heated on a hot plate to char the organic
material. The sample is heated in a muffle furnace at 550°C for several
hours. After cooling to room temperature the residue is dissolved by adding 2
mL of 1:1 HNO3 and evaporated to dryness. The residue is redissolved in 10
mL of 1:9 HNO3, filtered, and diluted to 50 mL in a volumetric flask. The
following data were obtained during a typical analysis for the concentration
of Na in a 4.0264-g sample of oat bran.
34
Determine the parts per million Na in the sample of oat bran.

2-5) Gluodenis describes the use of ICP to analyze samples containing Pb and Ni
in brass. The analysis for Pb uses external standards prepared from brass
samples containing known amounts of lead. Results are shown in the
following table.
What is the %w/w Pb in a sample of brass that gives an emission intensity of

9.25 × 104?
2-6) Five standard solutions were prepared for measuring the lead concentration
in two solutions, A and B. The two solutions A and B contain the same
concentration of magnesium used as an internal standard. The following
data were obtained:
From the data, calculate the lead concentration (mg/L) in the two sample
solutions, A and B.
35
Chapter 3: ATOMIC ABSORPTION SPECTROSCOPY (AAS)

3.1. The appearance of atomic absorption spectrum
As mentioned in the previous chapter, matter is formed by atoms that are
the smallest particles remaining to have properties of element. An atom consists
of a nucleus and electrons moving in the space around the nucleus. Under normal
conditions, atoms do not absorb or emit energy through radiation. Such atoms are
existing at ground state. If an atom is in gaseous state of unbounded one that is
irradiated by a beam of light having the wavelengths equal to those of light
emitted by itself in its atomic emission spectrum, they could absorb energy from
the light and convert to excited state having higher energy. This process is called
energy absorption one of atom, and it is the characteristic property of a gaseous
atom. The specrum generated in this process is called atomic absorption
spectroscopy.
If denoting energy of light absorbed by atom is ΔE, we have
36
ΔE = (Em – Eo) = hγ (3.1)

or:
∆ E=h . c /¿ (3.2)
Here,
E 0 and E m are energies of atom at ground state and excited state,
respectively;
h is the Plank’s constant; γ is the frequency of the light; and l s the
wavelength of absorption line.
As a result, corresponding to each value of DE that an atom has been
absorbed, we have an absorption line with the wavelength of li , meaning that

atomic absorption spectrum is also lines spectrum like atomic emission one.
It is noted that an atom does not absorb all radiation emitted by itself in the
emission process. The absorption process only occur for sensitive absorption
spectrum lines called characteristic spectrum lines of atom. Therefore, for these
spectrum lines the absorption and emission are two opposite process (Figure 3.1).
Accoding to Equation 3.1, the value of energy of DE is positive for the emission
process, and negative for the absorption process. So depending upon specific
conditions of energy source used for atomization leads to the dominance of one of
these processes.
In a mearurement of atomic absorption spectroscopy, a vapor of unbounded
gaseous atoms is initially formed, and then the cluster is irradiated by a beam of
light having the wavelengths equal to those of light emited in emission process of
the analyting elements. At that time the atom absorbs energy and generate its
apsorption spectrum.
37
Figure 3.1: AAS and AES processes of Na atoms

3.2. The intensity of absorption spectrum lines
Theoritical and experimental studies on the dependence of the intensity of
an absorption spectrum line of an element upon the concentration of an analyte in
a sample show that in a specific concentration range, the relationship between the
intensity and the concentration of this element in the vapor also complies with the
Lamber- Beer rule. It means that if the intensities of the beam of light before and
after passing through the vapor of atoms having the length of l and the
concentration of N are respectively denoted Io and I, we have:
I = Io.e –Kγ.N.l (3.3)
Where Kγ is atomic absorption coefficient of a spectrum line having the
frequency of γ.
Kγ is chraeristic for each absorption spectrum line of an element.
If the intensity of an atomic absorption spectrum line is indicated as D,
from (3.3), we have
D = log(Io/I) = 2,303.Kγ.N.l (3.4)
D is also called absporbance that depends upon the concentration N of the
atoms in the absorption medium and the length l of the medium. Due to the fact
that the length l in a AAS spectrometer is constant, the value of D now only
38
remain to depend upon concentration N of the atom in the absporption medium.

So Equation 3.4 can be rewritten as follows:
D = K.N (3.5)
Where K is an experimental coefficient whose value depends on
- The value of Kγ of the spectrum line;
- The temperature of absportion medium;
- The length of absorption medium.
Because in specific conditions the concentration N of atoms in absorption
medium is linearly proportional to the concentration C of analyte in sample, we
have
N = Ka.Cb (3.7)
Where Ka is an experimental constant depending upon all conditions
relating to the vaporization and atomization processes of sample; b is a constant
equal to 1 as the concentration of C is smaller than a specific value of Co
determined by experiments. In the case of the concentration of C is larger than the
value of Co, b is smaller than 1, meaning that N is not linearly proportional to C.
Combining two equations of (3.5) and (3.7), we have
D = k.Cb (3.8)
For the purpose of analyzing quantitatively, a measurement of AAS should
be conducted for samples having the concentration of analyte smaller than the
value of Co. So we have
D = kC
3.3. Principles and main components of an AAS measurement
As mentioned in the above section, an AAS analytical method is based on
the energy absorption of unbounded atoms of an element when irradiating a beam
of monochromatic light having a distinct wavelength through the gaseous atoms
of this element in absorption medium. So the principle of this method is based on
the following main stages:
1. Atomizing the sample:
In this the first stage, we need to chose an appropriate energy source and
experimental conditions to convert analyzing sample from initial state of solid or
liquid to gaseous state of unbounded atoms that are the absorption medium to
generate atomic absorption spectrum of the analyzing element. The component
used for the atomization of sample is called atomization system.
2. Generating the spectrum of atomic absorption
A beam of light emitted by the atomic emission of analyzing element is
irradiated through a vapor of the above formed atoms in which there are atoms of
39
analyzing element. These atoms would absorb some specific light and generate its
absportion spectrum. Here, parts of the initial intensity of the beam absorbed by
the atoms of analyzing element depend upon the concentration of these atoms in
the absorption medium. The emission beam of light of the analyzing element is
provided by a source as hollow cathode lamp (HCL).
3. Collimating, dispersing and choosing a disctint spectrum line of the
analyzing element to measure
The measured intensity is a absorption signal of the absorption spectrum
line. In a determinate limit of concentration, the ratio of the measured intensity
before and after passing through the absorption medium is linearly proportional to
the concentration C of the analyzing element in sample as presented in Equation
(3.8).
These above processes are principle of AAS measurement. So in order to
conduct an AAS measurement, an AAS spectrometer has to consist of the
following main components:
Component 1: A resonance emission source of the analyzing element
generates an emission beam of light to irradiate absorption medium containing a
vapor of the unbouned atoms of the analyzing element. For this purpose, a hollow
cathode lamp (HCL) is used.
A hollow-cathode lamp, in Figure 3.2, is filled with Ne (or Ar) at a pressure
of ≈ 1-5 Torr. The cathode is made of the element whose emission lines we want.
When are a potential of about 500V applied between the anode and the cathode,
gas is ionized and positive ions are accelerated toward the cathode. After
ionization occurs, the lamp is maintained at a constant current of 2–30 mA by a
lower voltage. Cations strike the cathode with enough energy to “sputter” metal
atoms from the cathode into the gas phase. Gaseous atoms excited by collisions
with high-energy electrons emit photons. This atomic radiation has the same
frequency absorbed by analyte in the flame or furnace. Atoms in the lamp are
cooler than atoms in a
40
Figure 3.2: A hollow cathode lamp (HCL)

Component 2: A system used to atomize the analyzing sample is called
sample atomization system for AAS. There are two common systems
corresponding to the technique of atomizing sample by flames (The AAS
measurement method using this technique is called F-AAS) and the technich of
atomizing sample by using a graphite furnace known as electrothermal
atomization (ETA-AAS or GF-AAS).
Component 3: it is a spectrometer including a monochomator used to
disperse a multi-chromatic beam into monochromatic lights and a photomultiflier
used to measure the intensity of a chosen beam of monochromatic light that
comes from the monochromator. The measured signals are amplified and treated
to obtain the values of absorbance of introduced sample.
Component 4: It is a system used to display the measured results of the
spectrum line. The output device could be a digital display used to process and
display the mesuared results, or a computer programmed to control the above
components. In brief, an F/GF-AAS spectrometer system can be illustrated in
Figure 3.4.
Digital display
or computer
Figure 3.4: Schematic diagram of a double beam F/GF-AAS

3.4. The technique of atomization by flames
The system used to atomize sample is a premix burner such as that in
Figure , in which fuel, oxidant, and sample are mixed before introduction into the
flame. Sample solution is drawn into the pneumatic nebulizer by the rapid flow of
oxidant (usually air) past the tip of the sample capillary. Liquid breaks into a fine
mist as it leaves the capillary. The spray is directed against a glass bead, upon
which the droplets break into smaller particles. The formation of small droplets is
termed nebulization. A fine suspension of liquid (or solid) particles in a gas is
called an aerosol. The nebulizer creates an aerosol from the liquid sample. The
mist, oxidant, and fuel flow past baffles that promote further mixing and block
large droplets of liquid. Excess liquid collects at the bottom of the spray chamber
and flows out to a drain. Aerosol reaching the flame contains only about 5% of the
initial sample.
41
The most common fuel-oxidizer combination is acetylene and air, which

produces a flame temperature of 2400–2700 K (Table 21-1). If a hotter flame is
required to atomize high-boiling elements (called refractory elements), acetylene
and nitrous oxide are usually used.
Figure 3.5: A premix burner

3.5. The technique of atomization by graphite furnace
An electrically heated graphite furnace is more sensitive than a flame and
requires less sample. From 1 to 100 µL of sample are injected into the furnace
through the hole at the center of Figure 3.6. Light from a hollow-cathode lamp
travels through windows at each end of the graphite tube. To prevent oxidation of
the graphite, Ar gas is passed over the furnace, and the maximum recommended
temperature is 2 550°C for not more than 7 s.
Compared with flames, furnaces require more operator skill to find proper
conditions for each type of sample. The furnace is heated in three or more steps to
properly atomize the sample. To measure Fe in the iron-storage protein ferritin, 10
µL of sample containing ~ 0.1 ppm Fe are injected into the furnace at ~ 90oC. The
furnace is programmed to dry the sample at 125 oC for 20s to remove solvent.
Drying is followed by 60s of charring at 1400 oC to destroy organic matter.
Charring is also called pyrolysis, which means decomposing with heat. Charring
creates smoke that would interfere with the Fe determination. After charring, the
sample is atomized at 2100oC for 10s. Absorbance reaches a maximum and then
decreases as Fe evaporates from the oven. The analytical signal is the time-
integrated absorbance during atomization. After atomization, the furnace is heated
42
to 2500oC for 3 s to clean out any remaining residue. The furnace is purged with
Ar or N2 during each step except atomization to remove volatile material. Gas
flow is halted during atomization to avoid blowing analyte out of the furnace.
When developing a method for a new kind of sample, it is important to record the
signal as a function of time, because signals are also observed from smoke during
charring and from the glow of the red-hot oven in the last part of atomization. A
skilled operator must interpret which signal is due to analyte so that the right peak
is integrated.
Figure 3.6: Diagram showing a cross-section of an electrothermal analyzer.
Figure 3.7: The graphite furnace and the graphite cuvette

3.6. Advantages and disavantages of AAS
Like other analytical methods, AAS methods also have not only many
advantages, but also some disadvantages. The advantages are:
- AAS methods, especially GF-AAS, have high analytical sensitivity and
selectivity. Due to the high analytical sensitivity, some of samples do not need to
concentrate analyte before measuring. Consequently, the analysis does not require
a large amount of sample, and it is not time-consuming. The sample
contamination is prevented due to the fact that we do not use a lot of chemicals for
the concentration.
43
One more advantage of AAS methods is that many elements can be

determined simultaneously in a sample with small analytical error.
Besides the above advantages, AAS methods have some disadvantages.
- These equipments are relatively expensive, especially GF-AAS system.
- Due to high analytical sensitivity, the sample contamination is very
significant for trace analysis. So chemicals used for AAS have to be analytically
pure.
-The main disadvantage of AAS is that it only reveals the elemental
composition of sample, but not shows us the bonding of these elements.
3.7. Object and applications of AAS
The main object of AAS methods is the trace analysis of metals in
inorganic and organic samples. With present equipments and techniques of AAS
methods, most of metals (about 65 elements) and some of nonmetals can be
determined quantitatively to the detection limit of ppb with the error smaller than
15%.
Besides metals, some metalloids such as Si, P, As, Se, Te are also
determined using AAS methods. Other nometals such as C, Cl, O, N are not
analyzed by these methods because the absorption spectrum lines of these
elements do not lie within the spectrum region that the monochromator of a
common AAS spectrometer can disperse, (190 - 900nm). For instance, the
absorption line of C is 165.70 nm, 134.70 nm for N, 130,20 nm for O; 134.78 for
Cl, and 180,70 nm for S.
Problems
3-1) Bonert and Pohl reported results for the atomic absorption analysis of several
metals in caustic suspensions produced during the manufacture of soda by
the ammonia-soda process. The concentration of Cu was determined by
acidifying a 200-mL sample of the caustic solution with 20 mL of
concentrated HNO3, adding 1 mL of 27% w/v H2O2 and boiling for 30 min.
The resulting solution was diluted to 500mL, filtered, and analyzed by flame
atomic absorption using matrix-matched standards. The results for a typical
analysis are shown in the following table.
44
Determine the concentration of Cu in the caustic suspension.

3-2) To measure Ca in breakfast cereal, 0.521 6 g of crushed Cheerios was ashed
in a crucible at 600°C in air for 2 h. The residue was dissolved in 6 M HCl,
quantitatively transferred to a volumetric flask, and diluted to 100.0 mL.
Then 5.00-mL aliquots were transferred to 50-mL volumetric flasks. Each
was treated with standard (containing 20.0 µg/mL Ca2+), diluted to volume
with H2O, and analyzed by flame atomic absorption. Construct a standard
addition graph and use the method of least squares to find wt% Ca in
Cheerios.
3-3) Quigley and Vernon report results for the determination of trace metals in sea
water using a graphite furnace atomic absorption spectrophotometer.
Calibration was achieved by the method of standard additions. The trace
metals were first separated from their complex, high-salt matrix by
coprecipitating with Fe3+. In a typical analysis a 5.00-mL portion of 2000-
ppm Fe3+ was added to 1.00 L of sea water. The pH was adjusted to 9 using
NH4OH, and the precipitate of Fe(OH) 3 allowed to stand overnight. After
isolating and rinsing the precipitate, the Fe(OH) 3 and coprecipitated metals
were dissolved in 2 mL of concentrated HNO 3 and diluted to volume in a 50-
mL volumetric flask. To analyze for Mn 2+, a 1.00-mL sample of this solution
was diluted to 100 mL in a volumetric flask. The following samples were
injected into the graphite furnace and analyzed
Report the parts per billion of Mn2+ in the sample of sea water.
45
3-4) An AAS method is employed for the determination of lead (Pb) in a sample
of adulterated paprika by the introduction of lead oxide (of the same colour).
An electrothermal atomic absorption instrument that provides a background
correction based upon the Zeeman effect is used. 0.01 g of the paprika
powder is placed in the tube of the graphite furnace. The determination of the
area peak absorbance was made at λ=283.3nm first in the absence and later
in the presence of a magnetic field. The value of the peak absorption
following background correction was 1220 (arbitrary units). Under the same
conditions, 0.01mL of a solution of 10 g/L Pb led to a value of 1000 in the
same units.
Calculate the % mass of lead in the sample of paprika under study.
3-5) Free cyanide in aqueous solution can be determined indirectly by atomic
absorption on the basis of its ability to dissolve silver as it passes through a
porous silver membrane filter at pH 12.
4Ag(s) + 8CN+ 2H2O + O2 → 4Ag(CN)2 + 4OH-
A series of silver standards gave a linear calibration curve in flame atomic
absorption with a slope of 807 meter units per ppm Ag in the standard. (The
meter units are arbitrary numbers proportional to absorbance, and ppm refers
to µg Ag/mL.) An unknown cyanide solution passed through the silver
membrane gave a meter reading of 198 units. Find the molarity of in the
unknown.
PART B: ELECTROCHEMICAL METHODS OF ANALYSIS

Electrochemical methods of analysis are based on the application of properties,
rules and electrochemical phenomena that relate to electrochemical reactions
46
occurring on the interface between electrode and analyte solution, and the
electrochemical properties of solution containing the cell.
These methods have developed for a lomg time, and especially in recent decades.
They are not only used in qualitative and quantitative analysis, but as a means to
investigate the theory of electrochemical processes and chemical reactions of
inorganic and organic substances.
Figure 0.1: Equippments for electrochemical methods of analysis

Electrochemical methods of analysis are equipped the following components:
- A cell containing a solution of analyte and electrolyte;
- Electrodes (at least two and sometimes three);
- Measuring device used to measure electric current or potential.
So all these methods have the same of theoretical bases on electrochemistry as
electrochemical equilibrium, electrodes, and electrode potential, and so on.
Categorizing electromethods of analysis
At present there are a lot of electrochemical methods of analysis, but they can
be divided into two main groups. A group is based on electrode processes, and the
other is based on measuring quantities that are not derived from electrode reaction
such as conductivity, resistance. In two these groups, the former is numerous and
diverse, and is divided into two subgroups.
+ The first subgroup includes electrochemical methods based on
electrochemical reactions occurring on electrodes under the condition of constant
electric current (nomarlly equal to zero).
+ The second subgroup includes electrochemical methods based on occurring
electrochemical reactions on electrodes in the presence of an electric current different
zero. This is a subgroup having electrolysis. This subgroup has many methods having
high sensitivity and being commomly applied in analysis.
Chapter 4: POTENTIOMETRIC METHODS
4.1. General principles
47
Analytical methods that are based upon potential measurements are termed
potentiometric methods. Potentialmetric methods are performed in a cell containing
analyte solution and two electrodes immersed in the solution including:
- 1 reference electrode has its potential constant (Eref = const)
- 1 indicator electrode indicates the potential of solution.
These two electrodes are connected to a potentiometer. The third component of
a potentiometric cell is a salt bridge that prevents the components of the analyte
solution from mixing with those of the reference electrode. The potential of the cell,
Ecell, is given by the equaton
E cell = ( Eind - E ref )
Here, Eind is the potential of indicator electrode;E ref is the potential of reference
electrode.
Because Eref is constant in during shows a schematic diagram of a cell for
potential measurement, the potential of potentiometric analysis.
the cell only depends upon the potential
of indicator electrode that contains the
information about the concentration of
the analyte in the measured solution. So
a potentiometric analysis involves
measuring a cell potential, correcting
this potential for the reference
potential, and correcting the anayte
concentration from the indicator
electrode potential.
It is noted that measuring
potential must be conducted in the
condition of electric current of zero or Figure 4.1: A schematic diagram of a
constant (I = 0, or = const). Figure 4.1 cell for potentiometric analysis
4.3. The equipment of potentiometry

The equipment of potentiometry consists of three main components as follows:
+ Electrochemical cell includes:
- a beaker containing the analyte solution,
- a indicator electrode
- a reference electrode
- a stirrer
+a potential meter
48
+an output display device

In present, the system could be equipped with an autosampler and a computer
so that the measuring process is automatically performed according to an available
programme.
4.3.1. Reference electrodes
The ideal reference electrode has a potential (versus the standard hydrogen
electrode) that is accurately known, constant, and completely insensitive to the
composition of the analyte solution. There are two reference electrodes have been
using in potentiometry including calomel reference electrodes and silver/silver
chloride reference electrodes.
4.3.1.1. Calomel reference electrodes
A calomel electrode can be represented
schematically as (Figure 4.2)
Hg/Hg2Cl2(s)│Hg2Cl2(saturated).KCl
(xM)║
Where x represents the molar
concentration of potassium chloride in
the solution. Three concentration of
potassium chloride are common, 0.1 M,
1 M and saturated (about 4.6 M) in
which the saturated solution is the most
widely used due to its easy preparation.
The potential of calomel electrode is
derived from the equilibrium
Hg2Cl2(s) + 2e-  2Hgo(l) + 2Cl- (aq)
Under the standard conditions,
the electrode potential of the saturated
Figure 4.2: Schematic diagram of a
calomel electrode (SCE) is 0.2444V saturated calomel electrode
(at 250C, p = 1atm, x = saturated).
4.3.1.2. Silver/silver chloride electrodes
A silver/silver chloride electrode consists of a silver electrode immersed in a
solution that is saturated both potassium chloride and silver chloride
Ag│AgCl(saturated),KCl (saturated)║
The potential of silver/silver choride is derived from the equilibrium
AgCl(s) + 1e-  Ago(s) + Cl-
49
Under the standard conditions (250C, p = 1atm, x = saturated), the potential of

this electrode is 0.199 V. A simple and easily constructed electrode is shown in Figure
4.3.
Figure 4.3: A schematic diagram of a silver/silver chloride electrode

Table 4.1. The potential of reference electrodes
Ecal (V) with x=

Temperature
(0C)
saturated 3.5 M 0.1M
15 0.2510 0.256 0.3370
20 0.2480 0.255 0.3359
25 0.2444 0.250 0.3356
30 0.2410 0.248 0.3351
35 0.2376 0.245 0.3344
EAg/AgCl (V) with x =

Temperarure
(0C)
Saturated 3.5 M 0.1M
15 0.2090 0.212 0.225
20 0.2040 0.208 0.221
50
25 0.1990 0.205 0.217
30 0.1940 0.201 0.212
35 0.1890 0.197 0.208
Table 4.1 shows the effect of the concentration of potassium chloride on the
potential of both the calomel electrode and silver/silver chloride electrode. In general,
the larger concentration and the higher temperature, the more decreased the potential
of electrode.
4.3.2. Indicator electrodes
An indicator electrode is one that there are such electrochemical processes
occurring on its surface as oxidation, reducrion, electron exchange, precipitation, and
so on. An ideal indicator electrode responds rapidly and reproducibly to changes in the
concentration of an analyte ion or group of analyte ions. Indicator electodes are of
three types: metallic, membrane, and ion-selective electrodes.
4.3.2.1. Metallic indicator electrodes
a) Electrodes of the first kind
An electrode of the first kind is a pure metal electrode that is in direct
equilibrium with its cation in the solution. This kind of electrode is used to keep track
of or determine the concentration of metal ions in solutions. For instance, the
equilibrium between metal Cu and its ion Cu2+ is:
Cu2+(aq) + 2e-  Cuo(s)

for which
o 0.0592
E ind = E Cu + log a Cu 2+
2
where Eind is the electrode potential of copper electrode and a Cu2+ is the activity
2+
of Cu ion.
b) Electrodes of the second kind
Metals not only serve as indicator electrodes for their own cations but also
respond to the activities of anions that form sparingly soluble precipitates or stable
complexes with such cations. For example, in silver/silver chloride electrode, the
electrode potential correlates to reproducibly with the activity of chloride ion in a
solution saturated with silver chloride. Here, the electrode reaction can be written as
AgCl (s) + 1e-  Ago(s) + Cl- có EoAgCl = 0.222 V
The Nernst expression for this process is
51
Eind = E oAgCl -
0.0592
1
(
log a Cl- )
This equation shows that the potential of a silver electrode is proportional to
the activity of chloride ion. Thus, in a solution saturated with silver chloride, a silver
electrod can serve as an indicator electrode of the second kind of chloride ion.
Another electrode of the second kind is mercury/mercury-EDTA complex
electrode. For example, when a small amount of HgY2- is added to a solution
containing Y4- (EDTA ion), the half-reaction at a mercury cathode is:
HgY2 + 2e-  2Hg(l) + Y4- có Eo = 0.21V
0.0592 a 4-
Eind = 0.21 - log Y
2 a HgY2-
for which
The stability constant for HgY2- is very large (6.3x1023), so the
concentration of the complex remains essentially constant over a large range of
Y4- concentration. The Nernst equation can be written as
0.0592
Eind = K - log a Y 4-
2
0.0592 1
K = 0.21 - log
2 a HgY 2-
Here,
This electrode is used to titrate EDTA.
c) Inert metallic electrodes for redox system
This kind of electrode is made from inert metal such as platinum, gold,
palladium responding to the potential of a redox system with which it is in contact.
For example, the equilibrium of a Ce4+ / Ce3+ system on platinum electrode is
Ce4+ + 1e-  Ce3+ Eo = 1.32V
And the potential of the platinum electrode is
0.059 a 4+
E ind = E o + log Ce
n a Ce3+
A platinum electrode is used as an indicator electrode for titration involving

standard cerium (IV) solution. For example, titrate (II) with Ce(IV) and vice versa
according to the following reaction:
Fe ( II ) + Ce ( IV ) = Fe ( III ) + Ce ( III )
4.3.2.2. Membrand electrodes

Typical electrode of this kind is the glass electrode for pH measurements. Until
now other types of membranes have developed to respond selectively to more than
52
two dozen other ions with high sensitivity. Such electrodes are also called ion sective
electrodes. In this section we will discuss some of them such as glass electrodes,
liquid membrand electrodes, crystalline-membrand electrodes. It is important to note
that there are fundermental differences both in design and in principle between
membrand electrodes and metal electrodes.
a) The glass electrode for pH measurements
The glass electrode used to measure pH is the most common ionselective
electrode. Figure shows diagram of a glass combination electrode for measuring pH
and a line diagram of this cell.
Figure 4.4: A glass combination electrode and a line diagram of this electrode
The pH-sensitive part of the electrode is the thin glass bulb or cone at the bottom
of the electrodes. The sample reference electrode at the left of the line diagram is the
Ag/AgCl electrode at the right of the combination electrode. The internal reference
electrode at the right side of the line diagram is the straight Ag/AgCl electrode at the
center of the electrode in Figure 4.4. The two reference electrodes measure the electric
potential difference across the glass membrane. The salt bridge in the line diagram is
the porous ceramic plug at the bottom-right side of the combination electrode
53
connecting analyte solution to the reference electrode. Two silver wires coated with
AgCl are visible inside the electrode.
- The composition and structure of glass membranes
The property and operation of glass membrane is a key factor for a glass pH
electrode. The composition of a glass membrane depends on the manufacturer. For
example, one of glass membranes, which has been widely used, consists of
aproximately 22% Na2O, 6% CaO and 72% SiO 2. This membrane shows excellent
specificity toward hydrogen ions up to a pH of about 9. Other glass formulations are
now in use in which sodium and calcium ions are replayced to various degrees by
barium and lithium ions. These membrane have superior selectivity and lifetime.
Figure 4.5 shows the irregular structure of the silicate lattice in glass. Negatively
charged oxygen atoms in glass can bind to cations of suitable size. Monovalent
cations, such as sodium and lithium can move sluggishly through the silicate lattice,
and are responsible for electric conduction within the membrane.
Both membrane surfaces must be hydrated before a glass electrode with
function as a pH electrode. Metal ions in these hydrated gel regions of the membrane
diffuse out of the glass and into solution. H+ can diffuse into the membrane to replace
the metal ions.
The reaction in which H+ replaces cations in the glass is an ion-exchange
equilibrium. A pH electrode responds selectively to H+ because H+ is the only ion that
binds significantly to the hydrated gel layer. The ion-exchange reaction can be written
as
H+(aq) + Na+(glass membrane)  Na+ (aq) + H+(surface membrane)
The potential difference between inner and outer silver-silver chloride electrodes
in Figure 4.3 depends on the chloride concentration in each electrode compartment
and on the potential difference across the glass membrane. Because the concentration
of Cl- is fixed in each compartment and because the concentration of hydrogen ions is
fixed on the inside of the glass membrane, the only variable is the pH of analyte
solution outside the glass membrane. So the potential difference E between two
reference electrodes can be expressed as a Nernst equation
Eb = constant + 0.05916loga1 = constant – 0.05916pH.
Here, a1 is the activity of H+ in the outer solution (analyte solution).
This equation states that the voltage of the ideal pH electrode changes by 59.16
mV for every pH-unit change of analyte activity at 25°C.
The ion- exchange process can be depicted in Figure 4.5,
54
Figure 4.5: Ion exchange at the interface between membrane and solution
- Chuẩn hóa một điện cực thủy tinh trong máy đo pH

Một điện cực thủy tinh pH nên được chuẩn hóa bằng hai (hoặc hơn) dung dịch
đệm chuẩn được chọn thích hợp sao cho pH của dung dịch đem đo nằm trong khoảng
pH của câc dung dịch chuẩn đó. Các dung dịch chuẩn này được pha để có các giá trị
pH ở một nhiệt độ xác định được cho trong bảng 2.2 với độ chính xác đến ±0.01 đơn
vị pH.
Before using a pH electrode, be sure that the air inlet near the upper end of the
electrode in Figure 4.4 is not capped. (This hole is capped during storage to prevent
evaporation of the reference electrode filling solution.) Wash the electrode with
distilled water and gently blot it dry with a tissue. Do not wipe it, because this action
might produce a static charge on the glass. To calibrate the electrode, dip it in a
standard buffer whose pH is near 7 and allow the electrode to equilibrate with stirring
for at least a minute. Following the manufacturer’s instructions, press a key that might
say “calibrate” or “read” on a microprocessor-controlled meter or adjust the reading of
an analog meter to indicate the pH of the standard buffer. Then wash the electrode
with water, blot it dry, and immerse it in a second standard whose pH is farther from 7
than the pH of the first standard. Enter the second buffer on the meter. Finally, dip the
electrode in the unknown, stir the liquid, allow the reading to stabilize, and read the
pH. Store a glass electrode in aqueous solution to prevent dehydration of the glass.
Ideally, the solution should be similar to that inside the reference compartment of the
electrode. If the electrode has dried, recondition it in dilute acid for several hours. If
the electrode is to be used above pH 9, soak it in a high-pH buffer.
If electrode response becomes sluggish or if an electrode cannot be calibrated
properly, try soaking it in 6 M HCl, followed by water. As a last resort, soak the
electrode in 20 wt% aqueous ammonium bifluoride, for 1 min in a plastic beaker. This
55
reagent dissolves glass and exposes fresh surface. Wash the electrode with water and
try calibrating it again. Avoid contact with ammonium bifluoride, NH 4HF2, which
produces a painful HF burn.
4.4. Potentiometric titration
A potentiometric titration involves measurement of the potential of a suitable
indicator electrode as a function of titrand volume. Potentiometric titrations provide
data that are more reliable than data from titrations that use chemical indicator and are
particularly useful with colored or turbid solution and for detecting the precence of
unsuspected species. Such titrations are also readily automated. Manual
potentiometric titrations of this kind, on the other hand, suffer from the disadvantage
of being more time-consuming than those involved indicators.
A typical apparatus for performing a manual potentiometric titration
between the analyte of Fe2+ and the titrant of Ce4+ is illustrated in Figure 4.6 . The cell
consists of two electrodes in which one is reference electrode and the other is indicator
electrode that could be inert metal such as Pt, Ag, Pd or glassy carbon or glass
membrand electrode.
Figure 4.6: The diagram of apparatus for potentiometric titration of Fe2+ with standard Ce4+
and the titration curve of 100.0 mL of 0.050M Fe 2+ in 1M HClO4 with 0.100M Ce4+ trong dung dịch
HClO4 1M.
Here, in the titration process the titrant is dropped from the buret to the beaker
containing the analyte leading to the change of its concentration. This change of
concentration results in the change of potential of the indicator electrode, especially
56
the abrupt change of potential around the equivalence point of the titration. This is the
signal that is used to detect the end point of the titration.
The potentiometric titration is applied for neutratization, precipitation, complex,
and redox reactions.
Problems
4-1) (a) Write the half-reactions for the silver-silver chloride and calomel reference
electrodes.
(b) Predict the voltage for the following cell.
4-2) Convert the following potentials. The Ag AgCl and calomel reference electrodes
are saturated with KCl.
(a) 0.523 V versus S.H.E = ? versus Ag│AgCl
(b) -0.111 V versus Ag │AgCl = ? versus S.H.E.
(c) -0.222 V versus S.C.E. = ? versus S.H.E.
(d) 0.023 V versus Ag │AgCl = ? versus S.C.E.
(e) -0.023 V versus S.C.E = ? versus Ag │AgCl.
4-3) For a silver-silver chloride electrode, the following potentials are observed:
E° = 0.222 V E(saturated KCl) = 0.197 V
Predict the value of E for a calomel electrode saturated with KCl, given that for
the calomel electrode is 0.268 V.
4-4) A cell was prepared by dipping a Cu wire and a saturated calomel electrode into
0.10 M CuSO4 solution. The Cu wire was attached to the positive terminal of a
potentiometer and the calomel electrode was attached to the negative terminal.
(a) Write a half-reaction for the Cu electrode.
(b) Write the Nernst equation for the Cu electrode.
(c) Calculate the cell voltage.
4-5) Explain why a silver electrode can be an indicator electrode for Ag + and for
halides.
4-6) A 10.0-mL solution of 0.050 0 M AgNO 3 was titrated with 0.025 0 M NaBr in the
cell
S.C.E. ║ titration solution │ Ag(s)
Find the cell voltage for 0.1, 10.0, 20.0, and 30.0 mL of titrant.
57
4-7) A solution prepared by mixing 25.0 mL of 0.200 M KI with 25.0 mL of 0.200 M

NaCl was titrated with 0.100 M in the following cell:
S.C.E. ║ titration solution │ Ag(s)
Call the solubility products of AgI and AgCl k I and kCl respectively. The answers
to parts (a) and (b) should be expressions containing these constants.
(a) Calculate when 25.0 mL of titrant have been added.
(b) Calculate when 75.0 mL of titrant have been added.
(c) Write an expression showing how the cell voltage depends on [Ag+]
4-8) Knowing that the standard potential of an electrode Cd/Cd 2+ vs. SHE, has the
value 0.404V, what happens to this potential if the electrode is plunged into an
aqueous solution of 0.01M CdSO4?
4-9) Calculate the molar concentration for the underlined component in the following
cell if the cell potential is measured at +0.294 V
Ag | AgCl (sat’d) | NaCl (0.1 M) || KI (x M)|I2(s)|Pt.
4-10) The concentration of NO3– in a water sample is determined by a one-point
standard addition using an NO3– ion-selective electrode. A 25.00-mL sample is
placed in a beaker, and a potential of +0.102 V is measured. A 1.00-mL aliquot
of a 200.0 ppm standard solution of NO 3– is added, after which the potential is
found to be +0.089 V. Report the concentration of NO3– in parts per million.
Chapter 5: VOLTAMMETRIC METHODS

5.1. Introduction
Voltammetry is a collection of techniques in which the relation between current
and voltage is observed during electrochemical processes. The earliest voltammetric
technique is polarography, developed by Jaroslav Heyrovsky in the early 1920s, an
58
achievement for which he was awarded the Nobel Prize in Chemistry in 1959. The
detection limit of this method lies within 10-3 - 10-5 Mol/L.
Due to the existence of capacitor current, many species can not be determined at
concentration smaller than 10-5 Mol/l leading to the appearance of square wave
polarography and pulse polarography. These later methods minimize significantly
capacitor current when measuring Faradic current so that they have higher sensitivity.
However, the detection limit gains n.10-6 Mol/L for the majority of substances, and
n.10-7 Mol/L for some of them.
Together with the development of other analytical methods, the stripping
voltamperometry methods were introduced. These methods are based upon the
combination between the electrolysis for concentration and polarography to increase
sensitivity. Stripping is the most sensitive voltammetric technique, because analyte is
concentrated from a dilute solution. The longer the period of concentration, the more
sensitive is the analysis. The detection limit for many metallic ions in aqueous
solution can be lowered to 10-12 M.
5.2. Voltametric measurements
In order to perform a voltametric
measurement for analyzing substances,
we have to dissolve analyte in a
suitable solvent (usually in water) to
obtain a sample solution that is then
added a large excess of inert supporting
electrolyte, for example KCl. The
resulting solution is measured using a
voltammeter making use of a three-
electrode potentiostat, such as that Figure 5.1: Diagram of voltammeter making
shown in Figure 5.1. use of a three-electrode potentiostat
A three-electrode potentiostat consits of at least three components as follows:

+ A electrochemical cell contains the resulting solution in which three electrodes,
including one working electrode, one standard reference electrode and one counter
electrode (or auxillary electrode), are immersed. The auxiliary electrode is generally a
platinum wire, and the reference electrode is usually a SCE or a Ag/AgCl electrode.
For the working electrode we can choose among several different materials, including
mercury, platinum, gold, silver, and carbon.
+ Power supply provides a suitable potential dependent on voltammetric
technique applied working electrode.
59
+ Voltmeter and current meter are used to measure the applied potential and the
Faradaic current when a reducing or oxidation reaction taking place at the woking
electrode.
A graph of current versus working electrode potential for a measuring sample
being oxidized or reduced at the working electrode is called voltammogram.
5.3. Current in Voltammetry
5.3.1. Faradaic current
When we oxidize an analyte at the working electrode by applying a suitable
potential, the resulting electrons pass through the potentiostat to the auxiliary
electrode, reducing the solvent or some other component of the solution matrix. If we
reduce the analyte at the working electrode, the current flows from the auxiliary
electrode to the cathode. In either case, the current from redox reactions at the
working electrode and the auxiliary electrodes is called a faradaic current.
The magnitude of Faradaic current is affected by three modes of mass transport
including diffusion, convection and migration.
Diffusion occurs whenever the concentration of an ion or molecule at the surface
of the electrode is different from that in bulk solution. The region of solution over
which diffusion occurs is the diffusion layer.
Convection occurs when we mechanically mix the solution, carrying reactants
toward the electrode and removing products from the electrode. The most common
form of convection is stirring the solution with a stir bar. Other methods that have
been used include rotating the electrode and incorporating the electrode into a flow-
cell.
The final mode of mass transport is migration, which occurs when a charged
particle in solution is attracted to or repelled from an electrode that carries a surface
charge. If the electrode carries a positive charge, for example, an anion will move
toward the electrode and a cation will move toward the bulk solution. Unlike diffusion
and convection, migration only affects the mass transport of charged particles.
Among the three described phenomena only diffusion can be related to
the concentration of the electro-active species, called also depolarizer.
The movement of material to and from the electrode surface is a complex
function of all three modes of mass transport. In the limit where diffusion is the only
significant form of mass transport, the current in a voltammetric cell is equal to
i=nFAD(Cbulk−Cx=0)/δ
Where n the number of electrons in the redox reaction, F is Faraday’s constant,
A is the area of the electrode, D is the diffusion coefficient for the species reacting at
60
the electrode, Cbulk and Cx=0 are its concentrations in bulk solution and at the electrode
surface, and δ is the thickness of the diffusion layer.
For the above equation to be valid, convection and migration must not interfere
with the formation of a diffusion layer.
We can eliminate migration by adding a high concentration of an inert
supporting electrolyte. Because ions of similar charge are equally attracted to or
repelled from the surface of the electrode, each has an equal probability of undergoing
migration. A large excess of an inert electrolyte ensures that few reactants or products
experience migration.
It is easy to eliminate convection by not stirring the solution.
5.3.2. Charging Current (capacitive current)
In addition to current resulting from redox reactions—what we call faradaic

current—the current in an electrochemical cell includes other, nonfaradaic sources.
Suppose the charge on an electrode is zero and that we suddenly change its potential
so that the electrode’s surface acquires a positive charge. Cations near the electrode’s
surface respond to this positive charge by migrating away from the electrode; anions,
on the other hand, migrate toward the electrode. This migration of ions occurs until
the electrode’s positive surface charge and the negative charge of the solution near the
electrode are equal. Because the movement of ions and the movement of electrons are
indistinguishable, the result is a small, short-lived nonfaradaic current that we call
the charging current (or capacitive current). Every time we change the electrode’s
potential, a transient charging current flows.
The capacitive current, is produced by the growth of a double electric
layer on the interface between the electrode and the solution. This double layer is
due to the high concentration of the supporting electrolyte in the solution and acts
as a condenser with high capacity. The total current flowing through the
electrode is finally due to the sum of the charging current (capacitive current)
of this condenser and the faradic current.
The capacitive current acts as a non specific background interference of
the faradic current, and sometimes can be higher than the latter, when the depolarizer
is present at low concentration in solution. In this cases the measure of the
faradic current is difficult and some electronic adjustment has to be used.
That’s why Polarography (and Voltammetry) is growth, as analytical technique,
only after the progress in the electronic field: so we can now affirm that the
development of this technique is strictly linked to the tentative to electronically
overcome problems due to capacitive current.
61
5.3.3. Residual Current
Even in the absence of analyte, a small, measurable current flows through an

electrochemical cell. This residual current has two components: a faradaic current
due to the oxidation or reduction of trace impurities and the charging current. Methods
for discriminating between the analyte’s faradaic current and the residual current are
discussed later in this chapter.
5.4. Voltammetric Techniques
5.3.1. DC polarography (Direct current polarography)
5.3.1.1. Principle of the method
Voltammetry conducted with a dropping-mercury electrode, is called
polarography. Polarography is an electrochemical method of analysis incorporating
features of electrolysis but distinct from it. In electrolysis the aim is to remove
completely a chosen constituent from the solution by passing an electric current
through it for a sufficient length of time. The electrodes have relatively large surfaces
and the solution is stirred to facilitate transport of electroactive material to the
electrode. In contrast to this, the electrolysis in Polarography is of short duration and
the electrode on which the constituents are placed out can be recovered virtually
unchanged.
`The electrochemical technique Polarography used in analytical chemistry,
involves measurements of current-voltage curves, obtained when is a dropping
mercury electrode or other micro electrode so that the currents are very small. Hence
the changes produced by Polarography are normally (with 5 to 20ml test solution) not
measurable, and the polarographic solution voltage is applied to electrodes immersed
in the solution being investigated.
62
Figure 5.2: A schematic diagram of DC polarography intrusment

One of the electrodes is an indicator electrode. It is a dropping mercury
electrode, consisting of a mercury drop hanging at the orifice of a fine bore glass
capillary. The capillary is connected to a mercury reservoir so that mercury flows
through the capillary at the rate of a few milligrams per second.
The life time of each drop is usually 3 to 5 seconds. Each drop forms a new
electrode; its surface is practically unaffected by processes taking place on the
previous drop. Hence, each drop represents a well-reproducible electrode with fresh
clean surface. The second electrode is a reference electrode; its potential remains
constant during the measurement. The potential at the indicator electrode varies in the
course of measurement of the current-voltage curve, because of the change of the
applied voltage.
In direct current polarography (DCP) a constant potential is applied during the
entire drop-life time. A current-voltage curve is constructed by applying a series of
potential steps, each step being synchronized with the drop fall. In most instruments,
63
however, linearly changing potential is applied, with a rate slow enough that the
change of potential throughout the drop-life time is about a few millivolts. The current
is measured at the end of the drop life.
Oxygen removal
Dissolved oxygen undergoes a two-step irreversible reduction at the dropping
electrode; the H202 produced in the first step is reduced to H20 in the second. The
two waves of equal size result, the first with a half-wave potential at about -0.14 V
and the second at about -0.9 V vs. S.C.E. The two half-reactions are somewhat slow.
As a consequence, the waves are drawn out over a considerable potential range. While
these polarographic waves are convenient for the determination of oxygen in
solutions, the presence of this element often interferes with the accurate determination
of other species. Thus, gas accomplishes this end. A stream of the same gas, usually
nitrogen, is passed over the surface during the analysis to prevent re-absorption.
5.3.1.2. Basis of quantitative and qualitative polarographic analysis
A polarogram can be plotted between the current flowing through the
polarographic cell against the increasing potential of the dropping electrode. A typical
polarogram is depicted in Figure 5.3.
The virtually flat upper portion of the polarographic wave is called its plateau
and the total current which flows through the cell at a potential on the plateau is called
the limiting current of the substance responsible for the wave.
The difference between the limiting current and residual is called the diffusion
current or wave height (id) of that substance and is a function of the concentration of
the electroactive material. The potential at which the current is one-half of the
diffusion current is called half-wave potential designated as E 1/2. The half-wave
potential of a reducible substance is independent of concentration and is the
characteristic of the nature of the reacting material.
This is the essential basis of quantitative and qualitative polarographic analysis.
A schematic diagram of a simple polarographic set up is shown in Fig. 5.2.
64
Figure 5.3: Typical polarogram obtained with the dropping mercury electrode.
Table 5.1: The half-wave potential of some inorganic and organic species
Metal Ions Matrix Half-wave potential (V)

Al (III) 0.1 M KCl - 1,700
Br - 0.1 M KNO3 - 0,121
Cd (II) 0.1 M KCl - 0,602
Cu (II) 0.1 M KCl - 0,340
Cu (II) 1M NaOH - 0,410
Fe (II) 0.1 M KCl - 1,300
Mn (II) 1M KCl - 2,365
Pb (II) 0.1 M KCl - 0,410
Zn (II) 0.1 M KCl - 0,995
65
The basis for the quantitative polarographic analysis is the existence of the linear
proportionality between the diffusion current and the concentration of the depolarizer
as given by Ilkovic equation.
Id = 607n . D1/2 . m 2/3 . t1/ 6 . C x
Where, Id = Average diffusion current, µA
n = Number of Faradays per mol involved in the electrode reaction
D = Diffusion coefficient of the electroactive material, (cm2/sec)
C = Concentration of the electroactive material, (millimoles/litre).
66
m = Flow rate of mercury through the capillary, (mg/sec).

τ = time between successive drops of mercury, (in seconds). I i là dòng khuyếch
tán giới hạn,
The diffusion current id is proportional to 'm2/3. τ 1/6, when the other factors in the
Ilkovic equation are constant. The drop time ‘τ’ and the flow rate ‘m’, which depend
upon the heights of mercury column, are constant due to the fact that the height of
mercury reservoir is hold constant. So we have
Ii = k . Cx
Here, k is a constant.
5.3.1.3. Advantagesof mercury electrode in DC polarography
The success of polarography was largely attributed to the fact that the dropping
mercury electrode has several advantages over solid metal such as platinum as an
indicator electrode. One of the most important of these is the high over potential for
the evolution of hydrogen on mercury. This makes it possible to study the reductions
of many substances including even the alkali and alkaline earth metal ions, which
could not possibly be deposited onto a platinum cathode without interference from the
simultaneous reduction of water or hydrogen ion. Unlike a solid electrode, a mercury
droplet has a perfectly smooth surface free from scratches or any other irregularities;
this makes the accurate calculation of the electrode area a rather simple matter.
Since, a fresh surface of the electrode is renewed regularly; absorbed or
deposited materials cannot accumulate on the electrode surface. Hence numerous
metals which are soluble in mercury electrode only negligibly small quantities of the
substance are deposited in the vicinity of dropping electrode.
For this reason there is no depletion of the depolarizer in the solution and
provided the volume of the electrolyzed solution is not too small, even after several
runs the curves obtained are identical. Finally, the dropping mercury electrode is much
less sensitive to mechanical disturbances than a stationary solid microelectrode
because a shock which would cause the latter to vibrate for sometime merely
dislodges one drop prematurely from a dropping electrode without having much effect
on the next. An even more serious defect of a stationary solid micro electrode is that
the concentration gradient resulting from the flow of current extends further and
further out into the
On the other hand, there are certain definite limitations associated with the
dropping electrode. The dropping mercury electrode may be applied over the range
+0.4 to about -2.0 volts with reference to S.C.E. At about +0.4 volt mercury dissolves
and gives an anodic wave; it begins to oxidize to mercury (I) ion. At potentials more
67
negative than about -1.8 Volts vs. S.C.E, visible hydrogen evolution occurs in acid
solutions and the usual supporting electrolytes commence to discharge.
5.3.1.4. Limitations of DC polarography
Ever since the development of DC polarography a great deal of efforts have been
undertaken to compensate several of the following three main limitations:
a) The drawback of direct current (DC) polarography as regards sensitivity is the
current required to change the potential of the dropping mercury electrode (DME) to
the required value: at a depolarizer concentration of 10 -5 mol.1-1, this time-averaged
charging current is comparable to the Faradaic current.
b) Lack of resolution of the sigmoid shaped waves (DC, NPP, etc.
polarographies). In order to sufficiently distinguish two polarographic waves so as to
be able to measure their corresponding half-wave potentials and limiting current
values, it is necessary that the former half-wave potentials should differ in more than
200 mV. All subsequent polarographic techniques, in which the obtained response is in
the form of a peak (derivative polarography, square-wave polarography, AC
polarography, differential polarography, differential pulse polarography, etc.) present a
higher resolution, being a difference of only about 50 mV sufficient to be able to
discriminate two substances presenting close half-wave potentials.
5.4. Pulse Polarography
In pulse methods the procedures are based on the application of pulse changes of
potential and the current response is measured at a suitable time relative to the time of
the pulse. The concept includes three methods: normal pulse polarography (NPP),
differential pulse polarography (DPP) and square-wave polarography (SWP). Pulse
techniques improve detection limits since they benefit from the different variation of
diffusion and capacitive current intensities with time: when carrying out
measurements at the pulse end, the capacitive current is practicably negligible, being
the value of the Faradaic current still significant.
5.4.1. Normal pulse polarography (NPP)
In normal pulse polarography (NPP) the mercury-drop electrode is held for
most of its duration at a constant potential Ein, at which no electrochemical reaction
takes place under given experimental conditions. The potential of interest E p is applied
in the last stage of the drop life, for a length of time t p (of the order of few
milliseconds). The values of Ein and tp are kept constant throughout the recording of
the polarogram and Ep is changed from drop to drop.
68
(a)
Figure 5.4: (a) Voltage profile of normal pulse polarography (NPP) (b) A typical NPP
polarogram
NPP keeps the DME at a constant potential before the start of the Faradaic current;
then, almost at the end of the drop life, a voltage pulse is applied, the amplitude of
which gradually increases from drop the drop. The pulse duration is about 200 ms.
The current is usually measured shortly and before the pulse end (about 20 ms after
the pulse edge) and is recorded versus the pulse amplitude. The current-potential
relationship for a reversible process is
Dx
Ii = n.F.S. . Cx
pt p
or: Ii = k1.Cx
Where S: The surface area of the dropping mercury electrode at the sampling
time (cm2)
F: Faraday’s constant
n: The electron number exchanged in the reaction at the eletrode
Dx: The diffusion coefficient of analyte
tp: The time length of a pulse (sec)
Cx: The concentration of analyte in the measuring solution
Ii : limitting current intensity
This mode is about seven times more sensitive than classical DC polarography.
The shape of the curve is that of a normal DC polarogram. The detection limit of NPP
is about 2·10-7 M.
5.4.2. Differential Pulse Polarography
This technique is comparable to normal pulse voltammetry in that the potential
is also scanned with a series of pulses. However, it differs from NPP because each
potential pulse is fixed, of small amplitude (10 to 100 mV), and is superimposed on a
slowly changing base potential. Current is measured at two points for each pulse, the
69
first point (1) just before the application of the pulse and the second (2) at the end of
the pulse. These sampling points are selected to allow for the decay of the nonfaradaic
(charging) current. The difference between current measurements at these points for
each pulse is determined and plotted against the base potential.
From analytical point of view, the sensitivity of DPP is even better than that of
NPP. The potential sequence on a single mercury drop and the potential sequence used
for recording an entire differential pulse polarogram are given in Fig.5.5. The current
is sampled twice during a drop life-time: (i) at t1, just before the pulse about 17 ms,
and (ii) at τ, just before the drop fall about 17 ms. The polarogram represents the
current difference as function of the base potential. The curve is
peaked shaped (Figure 5.6).
For Nernstian processes Ox + ne = Red and C*red = 0, the faradaic component
of the current at the peak Ii is
D ( s - 1) x C
Ii ( max ) = n.F.S . x
p.t p ( s + 1) x
k2
Where:
�nFΔE �
σ = exp � x �
�RT 2 �
DE : The magnitute of applied pulse
As a result, The height of the peak, Ii(max), is proportional to the
concentration of the electroactive species.
Ii ( max ) = k 2 .Cx
This equation is a base for the quantitative analysis of DPP method to
determine a lot of both inorganic and organic electroactive species with the sensitivity
up 10-8 Mol/l.
70
Figure 5.5 : Comparison of the voltage profile between NPP (left) and DPP (right)
NHCH3 NHCH3
N N
E1/2 = -0.265 V
+ 2H+ + 2e-
Cl N+ N
Cl
O-
C6H5 C6H5
E1/2 = -0.590 V + 2H+ + 2e-
H NHCH3 NHCH3
N N
E1/2 = -1.120 V
NH + 2H+ + 2e-
Cl Cl NH
C6H5 C6H5
Figure 5.6: Comparison of direct current (D.C) and different puled polarograph (DPP) of 1.2
x10-4 M chlodiazepoxide trong 3 ml of 0.05M H2SO4. Modulation amplitude = 50 mV.
5.5. Stripping voltammetry

5.5.1. Principle of the method
The concentration detection limits for classic d.c. polarography are only 10 -5 M
due to charging and other background currents. Therefore, this technique is limited in
use for quantitative concentration measurement. The discrimination against the
71
charging current in techniques such as normal pulse or differential pulse voltammetry

extends the detection limits down to 10-7-10-8M. However, the detecton limits of these
techniques can be lowered further (to 10 -10-10-11 M, i.e. sub parts per billion (ppb)) by
using as part of stripping voltametry.
Stripping voltammetry is an analytical technique that involves
(i) preconcentration (or deposition) step, in which analyte atoms/molecules are
deposited onto or into the electrode from the solution. Stirring the solution increases
the efficiency of this process. The concentration of the analyte on or in the electrode is
therefore is much higher than the concentration in the solution (up to 1000-fold
increase). The length of the preconcentration time determines the analyte
concentration, which in turn determines the detection limit of the technique.
(ii) Following a period of inactivity alowing homogeneous distribution of the
analyte and restoration of quiescent solution conditions, the analyte is then stripped
from the electrode by reversing the direction of the voltage sweep so that the species
is back into solution. Current measured during the stripping is proportional to the
quantity of analyte that was deposited.
The response due to a give analyte in the stripping step is proportional to the
concentration of that analyte in the solution. The conversion of the response to
concentration can be achived using a calibration curve, but the method of standard
additions is generally used.
In this method, a known amount of the analyte is added to the solution and the
experiment is repeated. For each sample, two or more standard additions are often
made.
5.5.2. Anodic stripping voltametry (ASV)
This is probably the most well-known stripping technique, and is used for
detection of metal cations, e.g., lead, cadmium, copper and zinc.
The preconcentration step involves the reduction of the Mn+ cations to the
metallic state at a mercury electrode, which leads to the formation of an M/Hg
amalgam.
Mn+ + ne-  Mo (M/Hg amalgam)
The M atoms are then reoxidized during the stripping step by scanning in a
positive (anodic) direction.
Mo (M/Hg amalgam) - ne-  Mn+ (aqueous solution)
The current during this step could be measured using DPP or square wave
polarography.
The peak potential in this scan is used for identification of the metal, and the
peak current is proportional to the concentration.
72
Figure 5.7 shows an anodic stripping

voltammogram of Cd, Pb, and Cu from
honey.
Figure 5.7: Anodic stripping voltammogram
of honey dissolved in water and acidified to
pH 1.2 with HCl. Cd, Pb, and Cu were
reduced from solution into a thin film of Hg
for 5 min at -1.4 V (versus S.C.E.) prior to
recording the voltammogram. [From Y. Li,
F.Wahdat, and R. Neeb, “Digestion-Free
Determination of Heavy Metals in Honey,”
Fresenius J. Anal. Chem. 1995, 351, 678.]
Two types of mercury electrode are commonly used, the Hanging Mercury
Drop Electrode (HMDE) and the Thin Mercury Film Electrode (TMFE).
 HMDE
Figure 5.8: Schematic diagram of a HDME (left) and a voltammetric instrument with HDME (right)
*TMFE
73
Figure 5.9: A diagram of a TMFE
This is often formed in the preconcentration step by the deposition of small

mercury droplets on a carbon glassy electrode.
There are advantages and disadvantages for each of these electrodes.
- The TMFE has a larger surface area-to-volume ratio than the HMDE, and this
leads to a higher analyte concentration ffor a given deposition time.
- The mercury drop is inherently more unstable than the mercury film.
- The HMDE gives reproducible results, one important requirement for any
analytical technique, better than the TMFE. Since a mercury film on a carbon surface
is in fact made up of the mercury droplets, the structure of the mercury layer is very
dependent on the nature of the surface of the carbon electrode, and it is not always
easy to obtain a reproducible surface for a solid electrode.
The importance of ASV has increased in recent years with the increased
awareness of the chronic toxic effects of metals and metal ions. Metals such as
cadmium, lead and mercury are toxic in trace amounts, and metals such as copper,
cobalt and zinc that are essential for living organisms in trace amounts are toxic at
higher concentrations. In addition, since metals are non-biodegradable and can not be
completely excreted, there is often accumulation in vital organs.
5.5.3. Cathodic stripping voltammetry (CSV)
The method of preconcentration for this technique is the formation of an
insoluble film on the surface of a mercury electrode. This film consists of a mercury
salt of the analyte anion An- (i.e., oxidation of the mercury electrode is required for
film formation). In the cathodic stripping step, the film is reduced, regenerating An-
(in solution) and mercury. Since the preconcentration occurs only on the surface of the
mercury electrode, either the HMDE or a mercury pool electrode is preferable to the
TMFE. In addition, the presence of mercury(II) ions in the solution (to generate the
TMFE) can lead to irreproducible results for CSV due to the formation of insoluble
74
mercury(II) salts in solution. CSV has been used for the detection of inorganic anion
such as halide ions, selenide and sulfide ions, and oxyanions such as [MoO4]2- and
[VO3]2-. Some organic compounds (e.g., thiols and nucleic acid bases also form
insoluble mercury(II) salts, and are therefore suitable for CSV.
5.5.4. Adsorptive stripping voltammetry
In the three previous techniques, preconcentration step for AdSV, no charge is
transfered, and accumulation is achieved via the inherent tendency of particular
molecules to adsorb to a mercury surface (e.g., due to hydrophobicity) ( however, it
should be noted that there is a potential dependence for the preconcentration step, and
the optimum preconcentration potential needs to be determined experimentally). This
technique has been promoted extensively for organic molecules (e.g., chlopromazine,
dopamine, bilirubin, triazine- containing pesticides and polychlorinated biphenyls. It
should be noted that the organic molecules listed are pharmaceuticals and
agrochemicals. The ability ot detect trace concentratons of such molecules has become
increasingly important in recent years, due to the increased efficacy of many
pharmaceuticals and the increased knowledge of the toxic effects of trace amounts of
many agrochemicals.
AdSV can also be used for detection of transition metal cations, through the
formation of insoluble complexes (e.g., dimethylglyoxime complexes of nickel(II) and
cobalt(II).
The stripping step in AdSV can be anodic or cathodic, so the direction of the
potential scan should be specified. However, this has led to some confusion in the
literature, since AdSV with a cathodic stripping step has sometimes been reffered to as
CSV. It is important to distinguish whether the preconcentration step is electrolytic
(oxidation of mercury followed by a mercury salt film formation in CSV) or non-
electrolytic (adsorption on the electrode surface in AdSV).
5.6. Amperometric titrations
5.6.1. Hydrodynamic voltammetry
Liner-scan voltammetry in which the solution or the indicator electrode is kept in
motion is called hydrodynamic voltammetry.
Hydrodynamic voltammetry is performed in several ways. One method involves
stirring the solution vigorously while it is in contact with a fixed microelectrode.
Alternatively, the microelectrode is rotated at a constant high speed in the solution.
As described in section , during an electrolysis, reactant is carried to the surface
of an electrode by three mechanism: (1) migration under the influence of an electric
field, (2) convection resulting from stirring or vibration, and (3) diffusion due to
concentration differences between the film of liquid at the electrode surface and the
75
bulk of the solution. In voltammetry, every effort is made to minimize the effect of
migration by introducing an excess of an inactive supporting electrolyte When the
concentration of supporting electrolyte exceeds that of the analyte by 50 to 100-fold,
the fraction of the total current carried by the analyte approaches zero.
As a result, the current at any point in the electrolysis is determined by a
combination of (1) the rate of mass transfort of the analyte A to the edge of the Nernst
diffusion layer by convection and (2) the rate of transport of A from the outer edge of
the diffusion layer to the electrode surface. Convection maitains a constant supply of
A at the outer edge of the diffusion layer, however. Thus, a steady state current that is
determined by the applied potential.
5.6.2. Amperometric titrations
Hydrodynamic voltammetry can be employed to estimate the equivalence point
of titrations, provided at least one of the participants or products of the reaction
involved is oxidized or reduced at a microelectrode. Here, the current at some fixed
potential in the limiting current region is measured as a function of the reagent
volume.
Figure 5.10b shows some typical forms of Amperometric titration curves.
Curves (a) represents a titration in which the analyte reacts at the electrode while the
reagent does not. Curve (b) is typical of a titration in which the reagent reacts at the
microelectrode and the analyte does not. Curve (c) corresponds to a titration in which
both the analyte and the titrant reacts at the microelectrode.
There are two types of amperometric electrode systems. One uses a single
polarizable microelectrode that is often a rotating platimum electrode coupled to a
reference electrode; the other uses a pair of identical solid state microelectrodes
immersed in a stirring solution.
A dropping mercury electrode may also be used, in which case the solution is not
stirred.
Figure 5.10a shows a typical cell arrangement for amperometric titrations with a
rotating platinum electrode.
Amperometric titrations with one indicator electrode have been confined to
titrations in which a precipitate or a stable complex is the product.
Precipitating reagents include silver nitrate for halide ions, lead nitrate for sulfate
ion, and several organic reagents such as 8-hydroxyquinoline dimethylglyoxime, and
cupferron for various metallic ions that are reducible at microelectrode.
Several metal ions have also been determined by titration with standard solutions
of EDTA.
76
The exception just noted involves titrations of some organic compounds such as
phenols, aromatic amine and olefins,; hydrazine; arsenic (III) and antimony(III) with
bromine.
A standard solution of potassium bromate is added slowly from buret to acidic
solution of the analyte that also contains an excess of potassium bromide. Bromine is
formed by the reaction
BrO3- + Br - + 6H + ��
� 3Br2 + 3H 2O
Before the equivalent point, the bromine reacts completly with the analyte, so no
current is observed. After this point a rapid increase in current takes place due to
electrochemical reduction at the microelectrode, which is applied a suitable potential
for the reduction of bromine,
Br2 + 2e-  2Br-
The use of a pair of identical metallic microelectrode to establish the equivalent
point in amperometric titrations offers the advantages of simplicity of equiment and
not having to prepare and maintain a reference electrode.
(a) (b)
77
Figure 5.10: (a) Experimental set up for amperometric titrations with a rotating patinum electrode;
(b) Typical amperometric titration curves.
5.7. Coulometric titration
Coulometric titration is a method of titration in which the titrating agent is
produced in a solution by electrolysis and the required amount of the agent is
determined by measuring the number of coulombs used in preparing it. A constant
current source called an amperostat is used to maintain a constant current during the
electrolysis.
For example, consider the coulometric titration of iron(II) at a platinum anod.
The titrating agent Ce4+ is generated by the electrolysis of Ce 3+ that is introduced at the
outset an excess amount:
Ce3+ →Ce4+ + e-
With stirring , the cerium(IV) is rapidly transported from the surface of the
electrode to the bulk of the solution, where it oxidizes an equivalent amount of
iron(II):
Ce4+ + Fe2+ →Ce3+ + Fe3+
If we have a means of determining when the reaction between Ce 4+ and Fe2+ is
complete, the concentration of Fe 2+ in the solution can be calculated based on the
number of coulombs used in preparing Ce 4+. Here, the number of coulombs (Q)
resulting from a constant current of I amperes operated for t seconds is
Q = I.t
5.7.2. Determining the end point
In general, indicators in volumetric titrations for detection of the end points can
be used in coulometric titrations as well.Thus, for the titration of iron(II) with
cerium(IV) as just described, an oxidation-reduction indicator such as 1.10-
phenantroline can be used.
For many specific coulometric titrations, the end point detection can be
performed by using a pair of identical platinum microelectrodes immersed in a stirring
solution.
Apparatus for coulometric titrations
Figure 5.11a presents shows a diagram of a coulometric titration apparatus that
includes a constant-current source, a titration vessel, a switch, an electric timer and a
current meter.
Cell for coulometric titrations
78
Figure 5.11: (a) Diagram of a coulometric titration (b) A coulometric titration cell
As can be seen from Figure 5.11B, a typicall coulometric titration cell consists of
a generator electrode at which the titrating reagent is generated and a counter
electrode to complete the circuit. The generator electrode is a platinum rectangle (or a
wire coil) having a relatively large surface area to minimize elects of polarization. If
the reaction products from the counter electrode interfere the titration, the counter
electrode is isolated from the reaction medium by a sintered glass disk. For example,
hydrogen gas formed at the counter electrode reacts rapidly with most oxidizing
reagents generated at the generator electrode. This interference leads to a positive
systematic error unless the gas is formed in a separate compartment.
5.7.3. Advantages of coulometric titrations
Compared with conventional titrations, coulometric titrations have some
advantages as follow:
- Eliminating the problems associated with the preparation,
standarization, and storage of standard solution.
- Titrating small amounts of analyte with high accuracy.
- A single constant current source provides titrating reagents for
precipitation, neutralization, comlex formation, or oxidation/reduction
titrations.
- Being easily automated.
Karl Fischer Titration of H2O
The Karl Fischer titration, which measures traces of water in transformer oil,
solvents, foods, polymers, and other substances, is developed by Karl Fischer, a
79
chemist working at a petrochemical company in Germany in the 1930’s. The titration

is usually performed by delivering titrant from an automated buret or by coulometric
generation of titrant. The volumetric procedure tends to be appropriate for larger
amounts of water (but can go as low as 1 mg H 2O) and the coulometric procedure
tends to be appropriate for smaller amounts of water.
5.8. The Karl Fischer titration
K-F titration involves two reactions:
•In the first reaction, an alcohol (usually methanol or ethanol), sulfur dioxide
(SO2) and a base (B) react to form an alkylsulfite intermediate:
CH3OH + SO2 + B  [BH]SO3CH3 (1)
•In the second reaction, the alkylsulfite reacts with iodine (I 2) and the water from
the sample:
[BH]SO3CH3 + I2 + H2O + 2 B  [BH]SO4CH3 + 2[BH]I (2)
•Since water and I2 are consumed in equimolar amounts in reaction 2, if you
know the amount of I2 consumed, you know the amount of water that was present in
the sample.
5.8.2. Types of K-F titration
• There are 2 types of K-F titration: Volumetric and Coulometric. Although the
endpoint of the reaction is marked by a persistence of the yellow (I 2) color, using the
eyes is not very accurate. Both methods use bipotentiometric titration to measure the
amount of I2 consumed by the water.
•Bipotentiometric titration is simply monitoring the extent of reaction by
measuring changes in electrical conductivity of the reaction solution.
• The difference between them is primarily in the way the I 2 is generated. While
in volumetric titration, the I2 is included with the reagents, the I 2 is generated at a
generator electrode in coulometric titration.
• Which one the customer chooses depends on the method they are following or
personal preferences (if any), the titration equipment they currently have in their lab,
and last but not least, the water levels in the sample, generally:
– Volumetric: 0.1 – 100% H2O
– Coulometric: 0.001 – 0.1% H2O
Coulometric or Volumetric
• While both techniques are based on the same two step reaction mechanism,
they differ in the way that Iodine is introduced to the reaction. The same K-F reactions
occur in Coulometric titration as do in Volumetric titration. However, in the
80
Coulometric system, the I2 is generated in situ rather than added as a reagent. It is

generally considered more sensitive to lower water levels
•The I2 is generated at the anode: 2 I- → I2 + 2e-
•The cathode reaction is: 2BH + 2e- → H2 + 2B
We illustrate the coulometric procedure in Figure 5.12, in which the main
compartment contains anode solution plus unknown. The smaller compartment at the
left has an internal Pt electrode immersed in cathode solution and an external Pt
electrode immersed in the anode solution of the main compartment. The two
compartments are separated by an ion-permeable membrane. Two Pt electrodes are
used for end-point detection.
Anode solution contains an alcohol, a base, SO2, I-, and possibly another organic
solvent. Methanol and diethylene glycol monomethyl ether
(CH3OCH2CH2OCH2CH2OH) are typical alcohols. Typical bases are imidazole and
diethanolamine. The organic solvent may contain chloroform, formamide, or other
solvents. The trend is to avoid chlorinated solvents because of their environmental
hazards. When analyzing nonpolar substances such as transformer oil, sufficient
solvent, such as chloroform, should be used to make the reaction homogeneous.
Otherwise, moisture trapped in oily emulsions is inaccessible. (An emulsion is a fine
suspension of liquid-phase droplets in another liquid.)
In a typical procedure, the main of H2O if the I2:H2O stoichiometry is
compartment in Figure 5.12 is filled 1:1.
with anode solution and the
coulometric generator is filled with
cathode solution that may contain
reagents designed to be reduced at the
cathode. Current is run until moisture
in the main compartment is consumed,
as indicated by the end-point detection
system described after the Example. An
unknown is injected through the
septum and the coulometer is run again
until moisture has been consumed. Two
moles of electrons correspond to 1 mol Figure 5.12: Apparatus for coulometric Karl
Fischer titration.
A bipotentiometric measurement is the most common way to detect the end point
of a Karl Fischer titration. The detector circuit maintains a constant current (usually 5
or 10 µA ) between the two detector electrodes at the right in Figure 5.12 while
81
measuring the voltage needed to sustain the current. Prior to the equivalence point, the
solution contains I-, but little I2 (which is consumed in Reaction 17-19 as fast as it is
generated). To maintain a current of 10 µA, the cathode potential must be negative
enough to reduce some component of the solvent system. At the equivalence point,
excess I2 suddenly appears and current can be carried at very low voltage by Reactions
A and B in Demonstration 17-2. The abrupt voltage drop marks the end point.
Anode: I3- + 2e- → 3I- Cathode: 3I- → I3- + 2e-
Problems
5-1) Explain the difference between charging current and faradaic current.
5-2) Explain what is done in anodic stripping voltammetry. Why is stripping the most
sensitive polarographic technique?
5-3) Write the chemical reactions that show that 1 mol of I 2 is required for 1 mol of
H2O in a Karl Fischer titration.
5-4) A dropping mercury electrode is regulated such that one drop falls every 4
seconds (t = 4s). The mass of the Hg corresponding to 20 drops is of 0.16 g.
Calculate the average mass flow of Hg in mg/s. If the flow is proportional to the
height of mercury in the reservoir, what would happen to the lifetime of the
drops if the height of Hg was increased three-fold?
5-5) The amount of sulfur in aromatic monomers can be determined by differential
pulse polarography. Standard solutions are prepared for analysis by dissolving
1.000 mL of the purified monomer in 25.00 mL of an electrolytic solvent, adding
a known amount of S, deaerating, and measuring the peak current. The following
results were obtained for a set of calibration standards
Analysis of a 1.000-mL sample, treated in the same manner as the standards,

gives a peak current of 1.77 µA. Report the amount of sulfur present in the
sample in milligrams per milliliter.
5-6) Differential pulse voltammetry at a carbon working electrode can be used to
determine the concentrations of ascorbic acid and caffeine in drug formulations.
In a typical analysis, a 0.9183-g tablet is crushed and ground into a fine powder.
82
A 0.5630-g sample of this powder is transferred to a 100-mL volumetric flask,

brought into solution, and diluted to volume. A 0.500-mL portion is then
transferred to a voltammetric cell containing 20.00 mL of a suitable supporting
electrolyte. The resulting voltammogram gives peak currents of 1.40 µA and
3.88 µA for ascorbic acid and caffeine, respectively. A 0.500-mL aliquot of a
standard solution containing 250.0 ppm ascorbic acid and 200.0 ppm caffeine is
then added. A voltammogram of this solution gives peak currents of 2.80 µA and
8.02 µA for ascorbic acid and caffeine, respectively. Report the number of
milligrams of ascorbic acid and of caffeine in the tablet.
5-7) Ratana-ohpas and co-workers describe a stripping analysis method for
determining levels of tin in canned fruit juices. Standards containing 50.0 ppb
Sn4+, 100.0 ppb Sn4+, and 150.0 ppb Sn4+ were analyzed, giving signals of 83.0,
171.6, and 260.2, respectively. A 2.00-mL sample of lychee juice was mixed
with 20.00 mL of 1:1 HCl/HNO 3. A 0.500-mL portion of this mixture was added
to 10 mL of 6 M HCl and the volume adjusted to 30.00 mL. Analysis of this
diluted sample gave a signal of 128.2. Report the parts per million of Sn 4+ in the
original sample of lychee juice.
5-8) Differential pulse polarography is used to determine the concentrations of lead,
thallium, and indium in a mixture. Since the polarographic peaks for lead and
thallium, and for thallium and indium overlap, a simultaneous analysis is
necessary. Peak currents (in arbitrary units) at potentials of –0.385 V, –0.455 V,
and –0.557 V were measured for a single standard solution and for a sample,
giving the results shown in the following table.
Report the concentrations of Pb2+, Tl+, and In3+ in the sample.

5-9) The concentration of Cu2+ in sea water may be determined by anodic stripping
voltammetry at a hanging mercury drop electrode. To determine the total amount
of Cu2+ in a sample of sea water, it is necessary to free any copper bound in
organic complexes. To a 20.00-mL sample of sea water is added 1 mL of 0.05 M
HNO3 and 1 mL of 0.1% H 2O2. The sample is irradiated with UV light for 8 h
and then diluted to volume in a 25-mL volumetric flask. Deposition of Cu 2+ takes
place at –0.3 V for 10 min, producing a peak current of 26.1 (arbitrary units). A
second 20.00-mL sample of the sea water is treated identically, except that 0.1
83
mL of a 5.00 µM solution of Cu2+ is added, producing a peak current of 38.4.

Report the concentration of Cu2+ in the sea water in parts per billion.
PART C: CHROMATOGRAPHY OF ANALYSIS

Chapter 6: FUNDERMENTALS OF CHROMATOGRAPHY
6.1. Introduction to chromatography
6.1.1 Definition of Chromatography
Tswett was the first to use the term "chromatography" derived from two Greek
words "Chroma" meaning color and "graphein" meaning to write.
. Tswett (1906) stated that " chromatography is a method in which the
components of a mixture are separated on adsorbent column in a flowing system”.
IUPAC definition (International Union of pure and applied Chemistry) (1993):
Chromatography is a physical method of separation in which the components to be
separated are distributed between two phases, one of which is stationary while the
other moves in a definite direction.
The stationary phase may be a solid, or a liquid supported on a solid or gel, the
mobile phase may be either a gas or a liquid.
6.1.2. Types of Chromatography
Different methods were attempted for classification of chromatography:
6.1.2.1. According to mechanism of separation
The mechanism of separation depends mainly on the nature of the stationary
phase. Based on separation mechanisms chromatography can be classified into:
 Adsorption Chromatography
It is the oldest and most common type of chromatography. The stationary phase
is a solid with adsorption power. Mixture components will be adsorbed on the surface
of the stationary phase with different powers and that account for separation. Silica gel
is the most common stationary phase in adsorption chromatography.
 Partition Chromatography
The stationary phase is a liquid forming a thin film on an inert solid acts as
support. The stationary liquid is usually more polar than the mobile liquid. The two
liquids must be immiscible with each other. Cellulose powder and wet silica gel are
examples of supports in partition chromatography that carry film of water act as
stationary phase. Partition chromatography is preferable over adsorption when dealing
with polar compounds.
 Ion Exchange Chromatography
84
It is used for separation of charged molecules. The stationary phase is an ion

exchange resin to which a cationic or anionic groups are covalently bonded. Ions of
opposite charges (counter ions) in the mobile phase will be attracted to the resin and
compete with the components of the mixture for the charged group on the resin. Both
the mixture components and the mobile phase must be changed. Mixture of Alkaloids
(compounds with positive charges) can be separated on anionic exchanger, while
mixture of organic acids (negative charges) can be separated using cationic exchanger.
Both types are used for desalination of water.
 Molecular Exclusion (Size Exclusion) Chromatography
Also called gel filtration or gel permeation chromatography, this technique
separates molecules by size, with the larger solutes passing through most quickly. In
the ideal case of molecular exclusion, there is no attractive interaction between the
stationary phase and the solute. Rather, the liquid or gaseous mobile phase passes
through a porous gel. The pores are small enough to exclude large solute molecules
but not small ones. Large molecules stream past without entering the pores. Small
molecules take longer to pass through the column because they enter the gel and
therefore must flow through a larger volume before leaving the column.
 Affinity Chromatography
It uses the affinity of proteins to specific ligands such as enzymes. The ligand is
attached to suitable polysaccharide polymer such as cellulose, agarose, dextran.
6.1.2.2. According to mobile phase
In this regard chromatography is classified into:
 Liquid Chromatography (LC)
The mobile phase is liquid. In case of separation by adsorption the stationary
phase is solid so it is called: Liquid-Solid Chromatography (LSC). If separation occurs
through partition the stationary phase is liquid so it is called: Liquid -Liquid
Chromatography (LLC).
 Gas Chromatography (GC)
Where the mobile phase is inert gas nitrogen or helium. Again if the stationary
phase is solid it is called: Gas–Solid Chromatography (GSC). When stationary phase
is liquid it is called Gas-Liquid Chromatography (GLC).
6.1.2.3. According to the technique (methods of holding the stationary phase)
 Planar or Plane Chromatography
In this type of chromatography the stationary phase is used in the form of layer.
Plane chromatography is further classified into:
a- Thin Layer Chromatography (TLC):
85
The stationary phase in the form of fine powder is spread on glass or plastic or
aluminum sheets.
b- Paper Chromatography (PC):
A specific type of papers is used as stationary phase in the form of sheets.
 Columnar or Column Chromatography (CC):
The stationary phase is held in to a tube made of glass or metal.
86
87
Figure 6.1: A diagram of a) column chromatography b) plannar chromatography
6.2. Generally theoretical fundamental of chromatography

6.2.1. Principle of chromatographic separation
Separation of two sample components in chromatography is based on their
different distribution between two non-miscible phases. The one, the stationary phase,
a liquid or solid, is fixed in the system. The other, the mobile phase, a fluid or a gas, is
streaming through the chromatographic system. In gas chromatography the mobile
phase is a gas, in liquid chromatography it is a liquid.
The molecules of the analytes are distributed between the mobile and the
stationary phase. When present in the stationary phase, they are retained, and are not
moving through the system. In contrast, they migrate with the velocity, v, of the
mobile phase when being there.
Due to the different distribution of the particular analytes the mean residence
time in the stationary phase differs, too, resulting in a different net migration velocity.
This is the principle of chromatographic separation.
Figure 6.2 illustrates how two components A and B are resolved on a column by
elution chromatography.
- A single portion of the samle, containing A and B, dissolved in the mobile
phase is introduced at the head of the column (at the time to) where components A and
B distribute themselves between two phases.
- Introduction of additional mobile phase (the eluent) forces the dissolved
portion of the sample down the column, where further patition between the mobile
phase and fresh portions of the stationary phase occurs (time t1). Partition between the
fresh eluent and the stationary phase takes place silmutaneously at the site of the
original sample.
- Furthur additions of the eluent carry solute molecules down the column in a
continuous series of transfers between the two phases.
- Because solute movement can occur only in the mobile phase, the average rate
at which a solute migrates depends on the fraction of time it spends in that phase.
- Component B is retained more strongly than A and the retention of component
A in the mobile phase is more likely than that of B leading to the fact that the average
rate through the column of A is larger than that of A.
- The resulting differences in rate cause the components in a mixture to separate
into bands along the length of the column.
88
- Passing a sufficient mobile phase through the column causes the isolation of
the separated species of A and B (time t2), and then each of them passes out the end of
the column where it can be detected by a suitable detector (time t3 and t4).
Chromatograms
- When the signal of a detector responding to solute concentration and placed at
the end of the column is plotted as a function of time, a series of symmetric peaks is
obtained as shown in Figure 6.2 b.
- A chromatogram is a graph showing the detector response as a function of
elution time.
- A chromatogram, is useful for both qualitative analysis, based upon the
position of a peak on the time axis, and quantitative analysis, based upon the area
under a peak.
Detector
Figure 6.2: (a) the separation of components A and B by column elution chromatography (b)
The output of the signal detector at the various stages of elution shown in (a)
6.2.2. Chromatographic characteristics of solute
6.2.2.1. Retention properties
 The chromatogram:
Solutes eluted from a chromatography column are observed with various
detectors described in later chapters. A chromatogram is a graph showing the
detector response as a function of elution time as presented in Figure 6.3.
89
Figure 6.3: A typical chromatogram
 Retention time, tr, for each component is the time needed after injection of the
mixture onto the column until that component reaches the detector.
 Retention volume Vr is the volume of mobile phase required to elute a
particular solute from the column.
 Retention volume: Vr = tr.uv
where uv is the volume flow rate (volume per unit time) of the mobile phase. The
retention volume Vr of a particular solute is constant over a range of flow rates.
Unretained mobile phase travels through the column in the minimum possible
time, designated tm. The adjusted retention time, t’r, for a solute is the additional time
required for solute to travel the length of the column beyond the time required by
unretained solvent:
t’r = tr - tm
In gas chromatography, tm is usually taken as the time needed for CH4 to travel
through the column (Figure ).
6.2.2.2. The capacity factor k’
For each peak in the chromatogram, the capacity factor, k’, is defined as
tr' tr - tm
k'= =
tm tm
The longer a component is retained by the column, the greater is the capacity factor. The capacity
factor in the above equation is equivalent to
90
where Cm is the concentration of solute in the stationary phase, Vm is the volume of the
stationary phase, Cs is the concentration of solute in the mobile phase, and Vs is the volume of the
mobile phase.
tr CV V
k'= = s s =K s
tm CmVm Vm
The quotient Cs/Cm is the ratio of concentrations of solute in the stationary and
mobile phases. If the column is run slowly enough to be near equilibrium, the quotient
Cs/Cm is the partition coefficient, K, introduced in connection with solvent extraction.
6.2.2.3. The selectivity factor (The relative retention)
tr' 2 k2' K 2
a= = =
tr' 1 k1' K1
 The selectivity factor (the relative retention) of two solutes is proportional to

the ratio of their partition coefficients. The greater the ratio of partition coefficients
between mobile and stationary phases, the greater the separation between two
components of a mixture. This relation is the physical basis of chromatography.
6. 3. Band broadening and column eficiency
6.3.1. Introduction
 Chromatographic efficiency is affected by the amount of band broadening that
occurs when a compound passes through the column.
 Solute moving through a chromatography column tends to spread into a
Gaussian shape with standard deviation σ (Figure 6.4).
 The longer a solute spends passing through a column, the broader the band
becomes. Common measures of breadth are (1) the width w 1/2 measured at a height
equal to half of the peak height and (2) the width w at the baseline between tangents
drawn to the steepest parts of the peak.
91
(a)
Figure 6.4: (a) the broadening of a peak (b) The width w at the baseline and w1/2 at the half of the
peak
6.3.2. Diffusion
 One main cause of band spreading is diffusion. The diffusion coefficient D
measures the rate at which a substance moves randomly from a region of high
concentration to a region of lower concentration.
 The number of moles crossing each square meter per second, called the flux, J,
is proportional to the concentration gradient:
mol dc
J ( 2 )=−D
m .s dx
Where:
 D is the diffusion coefficient,
 c is the concentration of the solute,
 x is the traveled distance.
 Diffusion in liquids is times slower than diffusion in gases. Macromolecules
such as ribonuclease and albumin diffuse 10 to 100 times slower than small
molecules.
 If solute begins its journey through a column in an infinitely sharp layer with m
moles per unit cross-sectional area of the column and spreads by diffusion as it
travels, then the Gaussian profile of the band is described by
m −x 2/( 4 Dt )
c= e
√ 4 π Dt
Where:
 c is concentration , (mol/m3)
92
 t is time, and x is the distance along the column from the

 current center of the band. (The band center is always in this equation.)
 From the above equation, the standard deviation of band is
σ =√ 2 Dt
6.3.3. Plate Heigh H
 Plate height is the constant of proportionality between the variance, σ2 , of the
band and
the distance it has traveled, x. The name came from the theory of distillation in
which separation could be performed in discrete stages called plates.
Plate height is also called the height equivalent to a theoretical plate.
H = σ2/L
Where L is the length of the packing column.
 Martin and Synge treated a chromatographic column as if it were made up of a

series of contiguous bubble-cap-like plates, within which equilibrium conditions
always prevail throughout a column during elution.
 This assumption can never valid in the dynamic state that exists in a
chromatographic column, where phase are moving past one another at such a pace
that sufficient time is not available for equilibrium.
 Plate height is approximately the length of column required for one
equilibration of solute between mobile and stationary phases.
 The smaller the plate height, the narrower the bandwidth.
 The ability of a column to separate components of a mixture is improved by
decreasing plate height. An efficient column has more theoretical plates than an
inefficient column. Different solutes passing through the same column have different
plate heights because they have different diffusion coefficients.
 Plate heights are ~0.1 to 1 mm in gas chromatography, ~ 10 µm in high-
performance liquid chromatography.
6.3.4. The number of theoretical plates N in a column
L L2
=
N= H σ2
L×W
σ=
Where L is the column length, and 4t R so:
2
16 t R
2
= ƯW
N
93
2
5 . 54×t R
2
or N = ƯW 1/ 2
Where:
 tR is the retention time of a peak,

 W is the width of the peak at its base (in units of time),
 W1/2 is the width of the peak at half-height.
This expression is used for experimental determination of the number of plates

in a column.
The larger the number of plates, or the smaller the plate height, the better the
separation efficiency of a column.
6.4. Variables that affect column efficiency
6.4.1. Effect of the mobile phase flow rate
 The magnitude of kinetic effects on column efficiency clearly depends on the

length of time the mobile phase is in contact with the stationary phase, which in turn
depends on the flow rate of the mobile phase.
 For this reason, efficiency studies have generally carried out by determining H
as a function of mobile phase velocity.
 The study results show a minimum in H at low flow rates for both gas and
liquid chromatographies, and in general, the flow rate at a minimum for liquid
chromatography is smaller that for gas chromatography.
6.4.2. Theory of band broadening
 Over the last years, an enormous amount of theoretical and experimental effort
has been devoted to developing quantitative relationships that account for the effects
of experimental variables such as linear velocity of mobile phase, diffusion coefficient
in mobile phase and in stationary phase, capacity factor, diameter of packing particles,
thickness of liquid coating on stationary phase, on column efficiency in GC and LC.
 The efficiency of most chromatographic columns can be approximated by the
Van Deemter equation:
B B
H = A+ + CS u x + CM u x = A + + Cu x
ux ux
Where:
 H is the plate height in centimeters;
 ux is the linear velocity of the mobile phase in centimeters per second;
94
 A is the multiple path term;

 B is the longitudinal diffusion coeficient;
 Cs and CM are mass-transfer coefficients for stationaryand mobile
phases, respectively.
Figure 6.5: The Van Deemter curve for a chromatographic column
6.5. Column resolution

6.5.1. Definition
 In Chromatography, the resolution Rs of a column is determined by the
following expression (Figure 6.6):
Δt R ΔV R t R 2−t R 1
= =
W av ư W av W 2 +W 1
Rs = 2
Where:
o ΔtR and ΔVR are the subtraction of two retention times or two retention
volumes of two peaks, respectively,
o Wav is the average width at the base of two peaks.
95
Figure 6.6: The separation efficiency corresponding to its resolution
 The resolution Rs of a column provides a quantitative measure of its ability to

separate two analytes. A resolution of 1.5 gives an essentially complete separation of
96
A and B ( the overlap is about 0.3%), where as a resolution of 0.75 does not and a
resolution of 1 gives the overlap between the two peaks is about 4%.
 The resolution for a given stationary phase can be improved by lengthening the
column, thus increasing the number of plates. However, this leads to an adverse
consequence of increasing the time required for the resolution.
6.6. Applications of chromatography
6.6.1. Qualitative analysis
Chromatography is widely used for recognizing the presence or absence of
components in mixture that contain a limited number of species whose indentities are
known. This can be carried out by comparing the retention time of the peak of
suspected species with that of its standard under the same chromatographic conditions,
and the comparison should be carried on columns having different stationary phases.
For example, as can be seen from Figure 6.7, peaks labeled 2,3,4,7 and 9 can be
identified as methyl, ethyl, n-propyl, n-butyl, and n-amyl alcohol, respectively, based
on their retention times.
Figure 6.7: Chromatographic qualitative analysis by comparing the retention time of an

unknown with that of its standard
On the other hand, because a chromatogram provides but a single piece of

information about each species in a mixture (the retention time), the application of the
technique in the qualitative analysis of complex samples of unknown composition is
limited.
This limitation has been largely overcome by linking chromatographic columns
directly with infrared and mass spectrometers. The resulting coupled instruments are
powerful tools for identifying the components of complex mixtures.
97
Another method to identify a peak is to compare its retention time with that of an
authentic sample of the suspected compound. The most reliable way to compare
retention times is by co-chromatography, in which an authentic compound is added to
the unknown. If the added compound is identical with a component of the unknown,
then the relative area of that one peak will increase. Identification is tentative only
when carried out with one column, but it is firmer when carried out on several
columns with different stationary phases.
It is important to note that while a chromatogram may not lead to positive
identification of the species in a sample, it often provides sure evidence of the absence
of species.
Thus, failure of a sample to produce a peak at the same retention time as a
standard obtained under identical conditions is strong evidence that the compound in
question is absent (or present at a concentration below the detection limit of the
procedure).
6.6.2. Quantitative analysis
Quantitative column chromatographv is based on a comparison of either the
height or the area of the analyte peak with that of one or more standards. For planar
chromatography, the area covered by the separated species serves as the analytical
variable. If conditions are properly controlled, these variables vary linearly with
concentration.
6.6.2.1. Internal normalization method
This method, also called ‘normalized to 100 per cent’ is used for mixtures for
which each component is producing a peak on the chromatogram, in order to be able
to make a complete assessment of the sample concerned. The solvent, if any, is
typically ignored.
Supposing that it is required to find themass concentrations of three compounds
1,2,3 in a mixture (Figure 6.8). The analysis is again carried out in two steps.
Step 1: Calculation of the relative response factors
A standard solution containing the three compounds 1,2, and 3 at known
concentrations C1, C2 and C3 is prepared. The chromatogram corresponding to the
injection of a volume V of this standard solution shows three peaks of area A1, A2 and
A3. These areas will be related to the masses m1, m2 and m3 of the compounds in
volume V.
98
Figure 6.8: Analysis by internal normalization method.

One of the compounds, 3 for example, is chosen as the substance for internal
normalization. This compound 3 will serve to calculate the relative response factors
K1/3 and K2/3 for compounds 1 and 2 with respect to 3. As previously deduced:
Given that mi=Ci.V, then the following expressions for K1/3 and K2/3 are obtained:
Step 2: Chromatogram of the sample – calculation of the concentrations

The next step consists to inject a sample of the mixture to be measured
containing
1,2 and 3. Labelling the elution peaks as A’1, A’2 and A’3 will gain direct access to
the percentage mass composition of the mixture represented by x1, x2 and x3 via
three expressions of the following form:
The condition of normalization being that: x1 +x2 +x3 =100.

For n components normalized to the component j:
6.6.2.2. External standard method

The first chromatogram is acquired from a standard solution (reference solution)
of known concentration Cref in a solvent. The standard and sample matrix should be as
similar as possible. A volume V of this solution is injected. Analysis conditions must
99
be identical. On the resulting chromatogram the area A ref of the corresponding peak is
measured.
The second chromatogram results from the injection of the same volume V of the
sample in solution, containing an unknown concentration of the compound to be
measured (conc. Cunk). The area of the corresponding peak is A unk. Since an identical
volume of both samples has been injected, the ratio of the areas is proportional to the
ratio of concen-
trations which depend upon the masses injected (mi = Ci. V).
Figure 6.9: Analysis by the external standard method.
The single point calibration method, as depicted in Figure 6.9, assumes that the
calibration line goes through the origin. Precision will be improved if the
concentrations of the reference solution and of the sample solution are similar.
Precision can be also improved if several injections of the sample and the
reference solutions are made, always using equal volumes. In a multilevel calibration,
equal volumes of a series of standard solutions are injected. This allows to get a
calibration curve of A = f (C), obtained by a regression method (linear least-square or
quadratic least-square). This leads to a more precise value for Cunk (Figure 6.10).
Figure 6.10: Quantitative analysis software for chromatography.
6.6.2.3. Internal standard method

100
For trace analysis it is preferable to use a method that relies on the relative
response factor of each compound to be measured against a marker introduced as a
reference. This means that any imprecision concerning the injected volumes, the
principal constraint of the previous method, is compensated.
Figure 6.11: Chromatograms for Internal standard method
The areas of the peaks to be quantified are compared with that of an internal
standard (designated by IS), introduced at a known concentration within the sample
solution.
Supposing that a sample contains two compounds 1 and 2 to be measured and
that compound (IS) represents the additional compound for use as an internal standard
(Figure 6.11).
Similarly, there are two steps for unknown calculation:
Step1: Calculation of the relative response factors
A solution containing compound 1 at known concentration C 1, compound 2 at
known concentration C2 and the internal standard IS at known concentration C IS is
prepared then injected onto the chromatograph. A1, A2, AIS will be the areas of the
elution peaks in the chromatogram due to the three compounds. If m 1, m2 and mIS
represent the real quantities introduced onto the column, of these three substances,
then three relations can be derived:
These ratios enable the calculation of the relative response factors of 1 and 2,
against IS and designated by K1/IS and K2/IS:
101
Since the injected masses mi are proportional to the corresponding mass

concentrations Ci (mi = Ci.V), the above equations can be rewritten as follows:
Step 2: Chromatogram of the sample – calculation of the concentrations
The second step of the analysis is to obtain a chromatogram for a given volume
of a solution containing the sample to quantify and to which has been added a known
quantity of internal standard IS. This will yield A’ 1, A’2 and A’IS. The areas of this new
chromatogram being obtained under the same operating conditions as previously. If
m’1, m’2 and m’IS represent the quantities of 1, 2 and IS introduced into the column,
then:
From the relative response factors calculated in the first experiment as well as
from the known concentration of the internal standard within the sample, C’ IS, this
leads to:
Expanding to n components it is easy to calculate the mass concentration of the

solute i using the following equation:
equally the percentage concentration of i can be expressed using equation:
This method becomes even more precise if several injections of the solution and
of the sample are carried out. Often a same volume of an IS stock solution is spiked
with all standards and samples.
In conclusion, this general and reproducible method demands nevertheless a
good choice of internal standard, which should have the following characteristics:
• it must be stable, pure and not exist in the initial sample
• it must be measurable, giving an elution peak well resolved on the
chromatogram
• its retention time must be close to that (or those) of the solute(s) to be
quantified
• its concentration must be close to, or above that of the analytes to quantify to
gain a linear response from the detector
• it must not interfere and co-elute with a sample component.
Problems
102
6-1) Explain why a component in a sample placed on a chromatography column

spreads out in a band as it moves through the column.
6-2) In chromatography, what is plate height? Is a larger or smaller plate height
desirable? Explain.
6.3) In a chromatographic analysis of low-molecular-weight acids, butyric acid elutes
with a retention time of 7.63 min. The column’s void time is 0.31 min. Calculate
the capacity factor for butyric acid. (Ans: 23.6)
6-4) In the same chromatographic analysis for low-molecular-weight acids considered
in 6-3, the retention time for isobutyric acid is 5.98 min. What is the selectivity
factor for isobutyric acid and butyric acid? (Ans: 1.29).
6-5) A chromatographic analysis for the chlorinated pesticide Dieldrin gives a peak
with a retention time of 8.68 min and a baseline width of 0.29 min. How many
theoretical plates are involved in this separation? Given that the column used in
this analysis is 2.0 meters long, what is the height of a theoretical plate?
(Ans:0.14 mm)
6-6) A mixture of benzene, toluene, and methane was injected into a gas
chromatograph. Methane gave a sharp spike in 42s, whereas benzene required
251 s and toluene was eluted in 333s. Find the adjusted retention time and
capacity factor for each solute. Also, find the relative retention of the two
solutes. (Ans: α = 1.39)
6-7) The following data were obtained for four compounds separated on a 20-m
capillary column.
Compound tr (min) w (min)
A 8.04 0.15
B 8.26 0.15
C 8.43 0.16
(a) Calculate the number of theoretical plates for each compound and the average
number of theoretical plates for the column.
(b) Calculate the average height of a theoretical plate.
(c) Explain why it is possible for each compound to have a different number of
theoretical plates.
6-8) Calculate the plate height (in mm) for a component placed on a 15-m column if it
has a retention time of 2.41 minutes and a baseline peak width of 0.069 min.
How many plates for the component are associated with this column?
6-9) Consider the following chromatography data for components separated on a 30-m
column.
103
a) Calculate the retention factor for component A.

b) Calculate the resolution for components A and B. Do these peaks overlap to a
significant extent? Explain.
6-10) A mixture placed in an Erlenmeyer flask comprises 6mL of silica gel and 40mL
of a solvent containing, in solution, 100mg of a non-volatile compound. After
stirring, the mixture was left to stand before a 10mL aliquot of the solution was
extracted and evaporated to dryness. The residue weighed 12mg. Calculate the
adsorption coefficient, K = CS/CM, of the compound in this experiment.
6-11) Moody studied the efficiency of a GC separation of 2-butanone on a dinonyl
phthalate column. Evaluating the plate height as a function of flow rate gave a
van Deemter equation for which A is 1.65 mm, B is 25.8 mm x mL x min –1, and
C is 0.0236 mm x min x mL–1.
(a) Prepare a graph of H versus u for flow rates in the range of 5–120 mL/min.
(b) For what range of flow rates does each term in the van Deemter equation
have the greatest effect?
(c) What are the optimum flow rate and the height of a theoretical plate at that
flow rate?
(d) For open tubular columns the A term is no longer needed. If the B and C terms
remain unchanged, what are the optimum flow rate and the height of a
theoretical plate at that flow rate?
(e) How many more theoretical plates will there be in the open tubular column
compared with the packed column?
6-12) A GC separation was conducted on a sample containing a pesticide analyte,
compound X. This sample was treated with an internal standard of Q, giving a
concentration of 15.0 ppm. A 1.0 µL injection onto the GC gave an FID respons
e of 1012 for Q and 3411 for X. A 1.0 µL standard solution of 30.0 ppm Q with
15.0 ppm of X was injected giving a response of 899 and 2791 respectively. Wh
at is the concentration of Q in the sample?
104
Chapter 7: GAS CHROMATOGRAPHY

7.1. Introduction
7.1.1. Overview
In gas chromatography, the components of a vaporized sample are separated as a
consequence of being partitioned between a mobile gaseous phase and a liquid or a
solid stationary phase held in a column.
There are two types of gas chromatography: gas-liquid chromatography (GLC)
and gas-solid chromatography (GSC).
In GLC the analyte is partitioned between a gaseous mobile phase and a liquid
phase immobilized on the surface of an inert solid packing or on the walls of a
capillary tubing. GSC is based on a solid stationary phase in which retention of
analytes occurs because of physical adsorption. While GLC finds widespread use in
all fields of science and its name is usually shortened to gas chromatography, the
application of OSC is limited because the nonlinear nature of the adsorption process
leading to severe tailing of elution peaks.
7. 1.2. Origins of Gas Chromatography
105
The development of GC as an analytical technique was pioneered by Martin and

Synge in 1941; they suggested the use of gas-liquid partition chromatograms for
analytical purposes.
The concept of gas chromatography was envisioned in the early forties but
unfortunately little notice was taken of the suggestion.
It was left to Martin himself and his co-worker A. T. James to bring the concept
to practical reality some years later in 1951, when they published their epic paper
describing the first gas chromatograph.
GC has developed into a sophisticated technique since the pioneering work of
Martin and James in 1951, and is capable of separating very complex mixtures of
volatile analytes.
In 1955 the first commercial apparatus for GLC appeared on the market.
Currently, nearly a million gas chromatographs are in use throughout the world.
7.2. Components of a GC installation
7.2.1. Introduction
A gas chromatograph is composed of several components within a special frame.
These components include the injector, the column and the detector, associated with a
thermostatically controlled oven that enables the column to attain high temperatures
(Figure 7.1). The mobile phase that transports the analytes through the column is a gas
referred to as the carrier gas. The carrier gas flow, which is precisely controlled,
enables reproducibility of the retention times.
Figure 7.1: Schematic diagram of a gas chromatograph
106
Analysis starts when a small quantity of sample is introduced as either liquid or

gas into the injector, which has the dual function of vaporizing the sample and mixing
it with the gaseous flow at the head of the column.
The column is usually a narrow-bore tube which coils around itself with a length
that can vary from 1 to over 100m, depending upon the type and the contents of the
stationary phase.
The column, which can serve for thousands of successive injections, is housed in
a thermostatically controlled oven. At the end of the column, the mobile phase (carrier
gas), passes through a detector before it exits to the atmosphere. Figure 7.2 shows a
chromatogram of analyzing a mixture of ketones by GC.
Figure 7.2: The GC chromatogram for a mixture of ketones
7.2.2. Carrier gas system

The mobilc-phase gas in GC is called the carrier gas and must be chemically
inert. Helium is the most common mobile-phase gas used, although argon, nitrogen,
and hydrogen are also used. These gases are available in pressurized tanks.
Pressure regulators, gauges, and flow meters are required to control the flow rate
of the gas. In addition, the carrier gas system often contains a molecular sieve to
remove impurities and water.
Flow rates are normally controlled by a two-stage pressure regulator at the gas
cylinder and some sort of pressure regulator or flow regulator mounted in the
chromatograph.
Inlet pressures usually range from 10 to 50 psi above room pressure, which lead
to flow rates of 25 to 150 mL/min with packed columns and 1 to 25 mL/min for open
tubular capillary columns. Generally, it is assumed that flow rates will be constant if
the inlet pressure remains constant.
107
Flow rates can be measured at the column head by using a soap-bubble meter.
Many modern computer-controlled gas chromatographs are equipped with electronic
flow meters that can be regulated to maintain the flow rate at the desired level.
Carrier gas
The nature of the carrier gas has no significant influence upon the values of the
partition coefficients K of the compounds between the stationary and mobile phases,
owing to an absence of interaction between the gas and solutes.
By contrast, the viscosity of the carrier gas and its flow rate have an effect on the
analytes’ dispersion in the stationary phase and on their diffusion in the mobile phase
(cf. Van Deemter’s equation), and by consequence upon the efficiency N and the
sensitivity of detection. Hydrogen is the carrier gas of choice. If this gas cannot be
used for safety reasons, helium may be substituted.
7.2.3. Sample Injection Systems
7.2.3.1. Sample introduction
The most common injection method is where a microsyringe is used (Figure 7.3)
to inject a very small quantity of sample in solution (e.g. 05L), through a rubber
septum into a flash vaporizer port at the head of the column.
Figure 7.3: Typical syringes used in GC
108
To better control the reproducibility of

the injections – simply changing the
user can lead substantial deviations,
when in manual mode – most
instruments are provides with
autosamplers in which the syringe
movements are automated (Figure 7.4).
The assembly operates in a cyclic
Figure 7.4: A typical autsampler fashion, taking the sample, injecting it
rapidly (0.2 s) and rinsing the syringe.
7.2.3.2.Injectors
The injector (Figure 7.5), which is the

sample’s entrance to the
chromatograph, has different functions.
Besides its role as an inlet for the
sample, it must vaporize, mix with the
carrier gas and bring about the sample
at the head of the column. The
characteristics of the injectors, as well
as the modes of injection, differ
according to column type.
The use of an automatic injection
Figure 7.5: A GC injector
system can significantly enhance
measurement precision.
7.3. Columns
There are two column types, which differ in their performance: packed columns
and capillary, also called open tubular, columns (Figure 7.6). For packed columns the
stationary phase is deposited or bonded by chemical reaction onto a porous support.
For capillary columns a thin layer of stationary phase is deposited onto, or bound to
the inner surface of the column.
109
In the past, the vast majority of gas chromatographic analyses used packed
columns. For most current applications, packed columns have been replaced by the
more efficient open tubular columns.
(b)
(a)
Figure 7.6: Typical capillary columns (a) and packed column (b)
Figure 7.7: The cross-section of (a) a packed column and (b) a capillary column
7.3.1. Packed columns

These columns, less commonly used today, have diameters of 1/8 or 1/4 inch
(3.18 and 6.35mm) and a length of between 1–3m (Figures 7.6b and 7.7a).
Manufactured from stainless steel or glass, the internal wall of the tube is treated to
avoid catalytic effects with the sample. They can withstand a carrier gas flow rate
within the range 10–40 mL/min. They contain an inert and stable porous support on
which the stationary phase can be impregnated or bounded (between 3 and 20 per
cent). This solid support, having a specific surface area of 2–8m2/g, is made of
spherical particles of around 0.2mm diameter, which are obtained from diatomites,
silicate fossils whose skeleton is chemically comparable to that of amorphous silica. A
110
linker assures the cohesion of the grains. After calcinations these diatomaceous earths
are often designated by the name of Chromosorb®. Other synthetic materials have
been developed such as Spherosil®, composed of tiny silica beads. All of these
supports have a chemical reactivity comparable to silica gel because of the presence of
silanol groups.
Although the performance of packed columns is more modest than capillary
columns, they are still usually employed for many routine analyses. Easy to
manufacture and with a large choice of stationary phases available, they are not
however, well adapted to trace analyses.
7.3.2 Capillary columns (open tubular columns)
The vast majority of analyses use long, narrow open tubular columns (Figure
7.6a) made of fused silica (SiO2) and coated with polyimide (a plastic capable of
withstanding 350°C) for support and protection from atmospheric moisture. Open
tubular columns offer higher resolution, shorter analysis time, and greater sensitivity
than packed columns, but have less sample capacity.
The wall-coated column in Figure 7.7b and 7.8a features a 0.1- to 5 µm thick
film of stationary liquid phase on the inner wall of the column. A support-coated
column has solid particles coated with stationary liquid phase and attached to the inner
wall. In the porous-layer column in Figure, solid particles are the active stationary
phase. With their high surface area, support-coated columns can handle larger samples
than can wall-coated columns. The performance of support-coated columns is
intermediate between those of wall-coated columns and packed columns.
(b)
Figure 7.8: (a) A wall coated column (b) Types of tubular columns
Column inner diameters are typically 0.10 to 0.53 mm and lengths are 15 to 100
m, with 30 m being common. Narrow columns provide higher resolution than wider
columns (Figure ) but require higher operating pressure and have less sample capacity.
Diameters of 0.32 mm or greater tend to overload the gas-handling system of a mass
spectrometer, so the gas stream must be split and only a fraction is sent to the mass
spectrometer. The number of theoretical plates, N, on a column is proportional to the
111
length. In Equation , resolution is proportional to and, therefore, to the square root of

column length (Figure ).
Thick films of stationary phase can shield analytes from the silica surface and
reduce tailing, but can also increase bleed (decomposition and evaporation) of the
stationary phase at elevated temperature. A thickness of 0.25 µm is standard, but
thicker films are used for volatile analytes.
Capillary columns made by deposition of these materials in the form of a fine
porous layer are called PLOT. They are employed to separate gaseous or highly
volatile samples. Columns containing graphitized carbon black have been developed
for the separation of N2, CO, CO2 and very light hydrocarbons. The efficiency of these
columns are very high.
7.3.3. Stationary Liquid phase
The choice of liquid stationary phase (Table 24-1) is based on the rule “like
dissolves like.” Nonpolar columns are best for nonpolar solutes (Table 24-2). Columns
of intermediate polarity are best for intermediate polarity solutes, and strongly polar
columns are best for strongly polar solutes. As a column ages, stationary phase bakes
off, surface silanol groups (Si-O-H) are exposed, and tailing increases. Exposure to
O2 from air at high temperatures also degrades the column. To reduce the tendency of
stationary phase to bleed from the column at high temperature, it is usually bonded
(covalently attached) to the silica surface or covalently cross-linked to itself.
Table 7.1: Some kinds of polysioxane
112
Squalane (nonpolar)
CH3 CH3 CH3 CH3
HC-(CH2)3-CH-(CH2)3-CH-(CH2)4-CH-(CH2)3-CH-(CH2)3-CH
CH3 CH3 CH3 CH3
The current phases correspond in principle to two families: the polysiloxanes

and the polyethylene glycols. Each category can be the object of minor structural
modifications.
Each of these phases can be used between a minimum temperature beneath
which concentration equilibria are too slow to occur, and a maximum temperature
above which degradation of the polymer occurs. This high limit depends on the film
thickness and the nature of the polymer.
- Polysiloxanes (also known as silicone oils and gums) are based upon a
repetitive backbone that consists of two hydrocarbon chains per silicon atom (Figure ).
There are about 20 different compositions of alkyl or aryl chains (methyl or phenyl) to
which can be incorporated further functional groups (e.g. cyanopropyl, trifluropropyl).
- Polyethyleneglycols (PEG): The best known representative of this family is
Carbowax® (Table 7.1). These polar polymers (Mr =1500 to 20 000 – for the
113
Carbowax 20M) can be used for deposition, impregnation or as bonded phases

(40<T<240/260 oC, depending on column diameter and film thickness).
Table 7.2: The polarity of compounds
7.4. Detector
7.4.1. Thermal conductivity detector
In the past, thermal conductivity
Table 7.3: Thermal conductivity of some gases detectors were most common in gas
chromatography because they are
simple and universal: They respond to
all analytes. Unfortunately, thermal
conductivity is not sensitive enough to
detect minute quantities of analyte
eluted from open tubular columns
smaller than 0.53 mm in diameter.
Thermal conductivity detectors are still
used for 0.53-mm columns and for
packed columns.
114
Figure 7.9: TCD detector
Thermal conductivity measures the ability of a substance to transport heat from a

hot region to a cold region (Table 7.3). Helium is the carrier gas commonly used with
a thermal conductivity detector. Helium has the second highest thermal conductivity
(after H2), so any analyte mixed with helium lowers the conductivity of the gas
stream. In Figure 7.9, eluate from the chromatography column flows over a hot
tungsten-rhenium filament. When analyte emerges from the column, the conductivity
of the gas stream decreases, the filament gets hotter, its electrical resistance increases,
and the voltage across the filament changes. The detector measures the change in
voltage. It is common to split the carrier gas into two streams, sending part through
the analytical column and part through a matched reference column. Each stream is
passed over a different filament or alternated over a single filament. The resistance of
the sample filament is measured with respect to that of the reference filament. The
reference column minimizes flow differences when temperature is changed.
Sensitivity increases with the square of the filament current. However, the maximum
recommended current should not be exceeded, to avoid burning out the filament. The
filament should be off when carrier gas is not flowing.
7.4.2. Flame Ionization Detector
In the flame ionization detector in Figure 7.10, eluate is burned in a mixture of
H2 and air. Carbon atoms (except carbonyl and carboxyl carbons) produce CH
radicals, which are thought to produce CHO+ ions and electrons in the flame.
CH + O →CHO+ + e-
Only about 1 in 105 carbon atoms produces an ion, but ion production is
proportional to the number of susceptible carbon atoms entering the flame. In the
absence of analyte, ~ 10-14 A flows between the flame tip and the collector, which is
held at +200 to 300 V with respect to the flame tip. Eluted analytes produce a current
115
of ~ 10-12 A, which is converted to voltage, amplified, filtered to remove high-

frequency noise, and finally converted to a digital signal.
Response to organic compounds is proportional to solute mass over seven orders
of magnitude. The detection limit is ~ 100 times smaller than that of the thermal
conductivity detector and is reduced by 50% when N2 carrier gas is used instead of
He. For open tubular columns, N2 makeup gas is added to the H2 or He eluate before it
enters the detector. The flame ionization detector is sensitive enough for narrow-bore
columns. It responds to most hydrocarbons and is insensitive to nonhydrocarbons such
as H2, He, N2, O2, CO, CO2, H2O, NH3, NO, H2S, and SiF4.
(a) (b)
Figure 7.10: (a) the cross-section and (b) the diagram of FID
7.4.3. Electron capture detector (ECD)

groups can react with thermal electrons
present to form negatively charged
ions.
The loss of such electrons is related to
the quantity of analyte in the sample.
In order to produce capturable (low
energy) thermal electrons, the carrier
gas is ionized by beta particles from a
Figure 7.11: A diagram of ECD radioactive source such as radioactive
The ECD is based on the phenomenon tritium or 63Ni in the cell. This electron
that organic molecules with flow produces a current, which is
electronegative functional groups, such collected and measured. When the
as halogens, phosphorous and nitro sample molecule is introduced into the
116
cell, electrons which would otherwise electrophile sample components. This

be captured at the electrode are change is what is recorded and
captured by the sample, resulting in measured for the chromatogram.
decreased current. . In comparison to a
signal without sample compounds, the
reduction in electron flow is
proportional to the quantity of
7.4.4. Mass Detector
A mass spectrometer is a powerful detector for both quantitative and qualitative
analysis of analytes in gas or liquid chromatography. In mass spectrometry, gaseous
molecules are ionized (usually to make cations) accelerated by an electric field, and
then separated according to their mass. The ionization process rts enough energy to
the molecule to break into a variety of fragments. A mass spectrum is a graph showing
the relative abundance of each fragement striking the detector of the mass
spectrometer. The horizontal axis of a mass spectrum is m/z, where m is the mass of
the ion and z is the number of charge it carries). Most ion carry one charge, so m/z is
usually equivalent to the mass of the ion in atomic units.
Eluate from the chromatograph is passed directly into a mass spectrometer,
which records the current from all ions of all masses over a wide, selected range.
A peak in the chromatogram is indentified by recording its mass spectrum as it
is eluted. Figure 7.12 shows the spectrum from peak 1 and identifies it as starting
material by comparison with a spectral library of known compounds. Mass of 1-
butanol is 74 atomic mass unit, and there is a tiny peak at this mass number in the
figure. This ion at m/z, which has not broken into any fragments is called the
molecular ion.
Figure 7.12 shows that peak 2 is identified as 1-bromobutane by noting that
natural bromine consisits of 50.52 79Br and 49.48% 81Br. The peak at m/z = 136 is the
unfragmented molecular ion, M+, with the formula C4H979Br+. The peak at m/z = 138
with slightly less intensity, is C4H981Br+.
117
Figure 7.12: The chromatogram and mass spectrum of 1-butanol and 1-

brombutan
The detector used widely to connect to an open tubular chromatography column

to record a spectrum of each peak as it is eluted is quadrupole mass spectrometer as be
seen in Figure 7.13. Species exiting the chromatography column pass through a heated
connector into the electron ionization source, which is pumped to maintain a pressure
of bar by using a high-speed turbomolecular or oil diffusion pump. Ions are
accelerated through a potential of 5–15 V before entering the quadrupole filter.
The quadrupole is one of the most common mass separators because of its low
cost. It consists of four parallel metal rods to which are applied both a constant voltage
and a radio-frequency oscillating voltage. The electric field deflects ions in complex
trajectories as they migrate from the ionization chamber toward the detector, allowing
only ions with one particular mass-to-charge ratio to reach the detector. Other ions
(nonresonant ions) collide with the rods and are lost before they reach the detector.
Rapidly varying voltages select ions of different masses to reach the detector.
Transmission quadrupoles can record 2–8 spectra per second, covering a range as high
as 4 000 m/z units.
118
Figure 7.13: Schematic illustration of a quadrupole mass detector
Problems
7-1) Why do capillary columns predominate in analytical GC?
7-2) What is temperature programming in GC? How does it gain an advantage over
single T separations?
7-3) What is the electron capture detector? Explain its basis for operation, why is N2
necessary? What types of species are detected with the ECD?
7-4) How does the capillary column configuration achieve its advantages over the
packed column setup in gas chromatography?
7- 5) What is an FID and how does it work?
What types of analytes does the FID respond to?
7-6) The analysis of the trihalomethanes CHCl 3, CHBr3, CHCl2Br, and CHClBr2 in
drinking water uses a packed column with a nonpolar stationary phase. Predict
the order in which these four trihalomethanes will elute.
7-7) A mixture of n-heptane, tetrahydrofuran, 2-butanone, and n-propanol elutes in
this order when using a polar stationary phase such as Carbowax. The elution
order is exactly the opposite when using a nonpolar stationary phase such as
polydimethyl siloxane. Explain the order of elution in each case.
7-8) Zhou and colleagues determined the %w/w H 2O in methanol by GC, using a
capillary column coated with a nonpolar stationary phase and a thermal
conductivity detector. A series of calibration standards gave the following
results.
119
(a) What is the %w/w H2O in a sample giving a peak height of 8.63?
(b) The %w/w H2O in a freeze-dried antibiotic is determined in the following
manner. A 0.175-g sample is placed in a vial along with 4.489 g of methanol.
Water in the vial extracts into the methanol. Analysis of the sample gave a peak
height of 13.66. What is the %w/w H2O in the antibiotic?
7-9) A GC‐FID analysis was conducted on a soil sample containing pollutant X. The
following separations were conducted:
tr(minutes)
peak area
Injection 1: 21.1 ppm Toluene Internal Standard 10.11 36,242
33.4 ppm X 14.82
45,997
Injection 2: 21.1 ppm Toluene Internal Standard 10.05
38,774
unknown concentration X 14.77
39,115
What is the concentration of X in the sample?
7-10) The amount of camphor in an analgesic ointment can be determined by GC
using the method of internal standards. A standard sample was prepared by
placing 45.2 mg of camphor and 2.00 mL of a 6.00 mg/mL internal standard
solution of terpene hydrate in a 25-mL volumetric flask and diluting to volume
with CCl4. When an approximately 2-µL sample of the standard was injected, the
FID signals for the two components were measured (in arbitrary units) as 67.3
for camphor and 19.8 for terpene hydrate. A 53.6-mg sample of an analgesic
ointment was prepared for analysis by placing it in a 50-mL Erlenmeyer flask
along with 10 mL of CCl4. After heating to 50 °C in a water bath, the sample was
cooled to below room temperature and filtered. The residue was washed with
two 5-mL portions of CCl4, and the combined filtrates were collected in a 25-mL
volumetric flask. After adding 2.00 mL of the internal standard solution, the
contents of the flask were diluted to volume with CCl4. Analysis of an
120
approximately 2-µL sample gave FID signals of 13.5 for the terpene hydrate and
24.9 for the camphor. Report the %w/w camphor in the analgesic ointment.
Chapter 8: LIQUID CHROMATOGRAPHY

8.1. Intrduction
Early liquid chromatography (LC) was carried out in glass columns with
diameters of 10 to 50 mm. The columns were packed with 50· to 500-cm lengths of
solid particles coated with an adsorbed liquid that formed the stationary phase. To
ensure reasonable flow rates through this type of stationary phase, the particle size of
the solid was kept larger than 150 to 200 µm; even then, flow rates were at best a few
tenths of a milliliter per minute. Thus, separation times were long – often several
hours. Attempts to speed up this classic procedure by application of vacuum or
pressure were not effective because increases in flow rates tended to increase plate
heights beyond the minimum in the typical plate height versus flow rate curve (see
Figure 26.8a)' decreased efficiencies were the result. '
Early in the development of liquid chromatography, scientists realized that major
increases in column effciency could be achieved by decreasing the particle size of
packings. It was not until the late 1960s, however, that the technology for producing
121
and using packings with particle diameters as small as 3 to 10 µm was developed. This
technology required sophisticated instruments operating at high pressures, which
contrasted markedly with the simple glass columns of classic gravity-flow liquid
chromatography. The name high-performance liquid chromatography (HPLC) was
originally used to distinguish these newer procedures from the original gravity-flow
methods. Today, virtually all LC is done using pressurized flow, and we use the
abbreviations LC and HPLC interchangeably.
LC is the most widely used of all of the analytical separation techniques. The
reasons for the popularity of the method are its sensitivity. its ready adaptability to
accurate quantitative determinations, its ease of automation, its suitability for
separating nonvolatile species or thermally fragile ones, and above all, its widespread
applicability to substances that are important to industry, to many fields of science,
and to the public.
In several basic types of chtomatography the mobile phase is a liquid. These
types are often classified by separation mechanism or by the type of stationary phase.
The varieties include: (1) partition chtomatography (2) adsorption, or liquid-solid
chtomatography (3) ion exchange, or ion chtomatography; (4) size exclution
chromatography; (5) affinity chtomatography, and (6) chiral chromatography; most of
this chapter deals with column applicatiors of these important types of
chromatography. The final section, however, presents a brief description of planar
liquid chromatography because this technique provides a simple and inexpensive way
of determining likely optimal conditions for column separations.
8.2. Main components of a HPLC
Figure 8.1 illustrates main components of a currently common HPLC
chromatography.
122
Figure 8.1: Schematic of a modular HPLC instrument (model HP 1200)
8.2.1. Mobile-Phase Reservoirs and Solvent Treatment Systems

A modern LC apparatus is equipped with one or more glass reservoirs, each of
which contains 500 mL or more of a solvent. Provisions are often included to remove
dissolved gases and dust from the liquids. Dissolved gases can lead to irreproducible
flow rates and band spreading; in addition, both bubbles and dust interfere with the
performance of most detectors. Degassers may consist of a vacuum pumping system, a
distillation system, a device for heating and stirring, or as shown in Figure 8.1, a
system for sparging, in which the dissolved gases are swept out of solution by fine
bubbles of an inert gas that is not soluble in the mobile phase. Often the systems also
contain a means of filtering dust and particulate matter from the solvents to prevent
these particles from damaging the pumping or injection systems or clogging the
column. It is not necessary that the degassers and filters be integral parts of the HPLC
system as shown in Figure 8.1. For example, a convenient way of treating solvents
before introduction into the reservoir is to filter them through a millipore filter under
vacuum. This treatment removes gases as well as suspended matter.
An elution with a single solvent or solvent mixture of constant composition is
termed an isocratic elution. In gradient elution, two (and sometimes more) solvent
systems that differ significantly in polarity are used and varied in composition during
the separation. The ratio of the two solvents is varied in a preprogrammed way,
sometimes continuously and sometimes in a series of steps. Modern HPLC
instruments are often equipped with proportioning valves that introduce liquids from
two or more reservoirs at ratios that can be varied continuously (Figure 8.1). The
volume ratio of the solvents can be altered linearly or exponentially with time.
123
8.2.2. Pumping Systems

The requirements for liquid chromatographic pumps include (1) the generation
of pressures of up to 6000 psi, or 414 bar, (2) pulse-free output, (3) flow rates ranging
from 0.1 to 10 mUmin, (4) flow reproducibilities of 0.5% relative or better, and (S)
resistance to corrosion by a variety of solvents.
The quality of a pump for HPLC is measured by how steady and reproducible a
flow it can produce. A fluctuating flow rate can create detector noise that obscures
weak signals.
Figure 8.2 shows a pump with two sapphire pistons that produce a
programmable, constant flow rate up to 10 mL/min at pressures up to 40 MPa (400
bar, 6 000 pounds/inch2).
Solvent at the left passes through an electronic inlet valve synchronized with the
large piston and designed to minimize the formation of solvent vapor bubbles during
the intake stroke. The spring-loaded outlet valve maintains a constant outlet pressure,
and the damper further reduces pressure surges. Pressure surges from the first piston
are decreased in the damper that “breathes” against a constant outside pressure.
Pressure surges are typically "1% of the operating pressure. As the large piston draws
in liquid, the small piston propels liquid to the column. During the return stroke of the
small piston, the large piston delivers solvent into the expanding chamber of the small
piston. Part of the solvent fills the chamber while the remainder flows to the column.
Delivery rate is controlled by the stroke volumes.
Figure 8.2: High-pressure piston pump for HPLC [Courtesy Hewlett-Packard Co., Palo Alto, CA.].
Gradients made from up to four solvents are constructed by proportioning the

liquids through a four-way valve at low pressure and then pumping the mixture at high
pressure into the column. The gradient is electronically controlled and programmable
in 0.1 vol% increments.
8.2.3. Injection Valves
124
The injection valve in Figure 8.3 has interchangeable sample loops, each of
which
holds a fixed volume. Loops of different sizes hold volumes that range from 2 to 1000
µL. In the load position, a syringe is used to wash and load the loop with fresh sample
at atmospheric pressure. High-pressure flow from the pump to the column passes
through the segment of valve at the lower left. When the valve is rotated 60°
counterclockwise, the content of the sample loop is injected into the column at high
pressure.
Figure 8.3: Injection valve for HPLC. Replaceable sample loop comes in various fixed-volume sizes.
8.2.4. HPLC Column
The HPLC equipment in Figure 8.1 decreases the viscosity of the solvent,
uses steel or plastic columns that are 5– thereby reducing the required pressure
30 cm in length, with an inner diameter or permitting faster flow.
of 1–5 mm (Figure 8.4). Columns are The most common HPLC column
expensive and easily degraded by dust diameter is 4.6 mm. There is a trend
or particles in the sample or solvent and toward narrower columns (2 mm, 1
by irreversible adsorption of impurities mm, and capillary columns down to 25
from the sample or solvent. Therefore, µL) for several reasons. Narrow
the entrance to the main column is columns are more compatible with
protected by a short guard column mass spectrometers, which require low
containing the same stationary phase as solvent flow. Narrow columns require
the main column. Fine particles and less sample and produce less waste.
strongly adsorbed solutes are retained
in the guard column, which is
periodically replaced. Heating a
chromatography column usually
125
Figure 8.4: HPLC column with replaceable

guard column [Courtesy Upchurch Scientific,
Oak Harbor,WA.]
8.2.4. HPLC Detector

An ideal detector of any type is sensitive to low concentrations of every analyte,
provides linear response, and does not broaden the eluted peaks. It is also insensitive
to changes in temperature and solvent composition. To prevent peak broadening, the
detector volume should be less than 20% of the volume of the chromatographic band.
Gas bubbles in the detector create noise, so back pressure may be applied to the
detector to prevent bubble formation during depressurization of eluate.
8.2.4.1. Detector UV- VIS
common HPLC detector, because many
solutes absorb ultraviolet light. Simple
systems employ the intense 254-nm
emission of a mercury lamp. More
versatile instruments have deuterium,
xenon, or tungsten lamps and a
monochromator, so you can choose the
optimum ultraviolet or visible
Figure 8.5: A Flow cell
wavelength for your analytes.
An ultraviolet detector using a flow cell
such as that in Figure 8.5 is the most
126
8.2.4.2. Photodiode Aray Detector
 The system in Figure 8.6 is able to use a photodiode array to record the spectrum
of each solute as it is eluted. Dual-beam optical system uses grating polychromator, one
diode array for the sample spectrum, and another diode array for the reference spectrum.
By this means, it was possible to determine which compounds were in each peak.
In a photodiode array spectrophotometer (Figure 8.6), all wavelengths are recorded
simultaneously, allowing more rapid acquisition of the spectrum or higher signal-to-noise
ratio, or some combination of both.
Each diode receives a different wavelength, and all wavelengths are measured
simultaneously.
Figure 8.6: Photodiode array ultraviolet detector for HPLC. (a) Dual-beam optical system (b)
chromatogram of sample containing 0.2 ng of anthracene, with detection at 250 nm.
Full-scale absorbance is 0.001. (c) Spectrum of anthracene recorded as it emerged from
the column. [Courtesy Perkin-Elmer Corp., Norwalk, CT.]
8.2.4.3. Fluorescence detector

The simplest detectors use a mercury excitation source and one or more filters to
isolate a band of emitted radiation. More sophisticated instruments are based on a xenon
source and use a grating monochromator to isolate the fluorescence radiation. Laser-
induced fluorescence is also used because of its sensitivity and selectivity.
An inherent advantage of fluorescence methods is their high sensitivity, which is
typically greater by more than an order of magnitude than most absorption procedures.
This advantage has been exploited in LC for the separation and determination of the
components of samples that fluoresce.
8.2.4.4. Refrative index detector
Figure 8.7 shows a schematic diagram of a differential refractive-index detector in
which the solvent passes through one half of the cell on its way to the column; the eluate
then flows through the other chamber. The two compartments are separated by a glass
plate mounted at an angle such that bending of the incidcnt beam occurs if the two
solutions differ in refractive index. The resulting displacement of the beam with respect
to the photosensitive surface of a detector causes variation in the output signal, which,
when amplified and recorded, provides the chromatogram.
 Refractive-index detectors have the significant advantage of responding to nearly
all solutes. That is, they are general detectors analogous to flame or thermal conductivity
detectors in gas chromatography. In addition, they are reliable and unaffected by flow
rate. They are, however, highly temperature dependent and must be maintained at a
constant temperature to a few thousandths of a degree centigrade. Furthermore, they are
not as sensitive as many other types of detectors and are not compatible with gradient
elution methods. Refractive-index detectors find application in the determination of some
analytes such as sugars.
Figure 8.7: Deflection-type refractive index detector.
8.3. HPLC methods

8.3.1. Adsorption chromatography
Bare silica can be used as the stationary phase for adsorption chromatography. For
the majority of applications silica gel still represents the basic material used to pack
HPLC columns. It is a rigid, amorphous solid having a composition formula SiO 2.nH2O
(where n is very close to 0), quite different from natural crystalline silica (SiO 2), which is
only a distant precursor of it. Silica gel is in the form of spherical particles, sometimes
porous, with a diameter of between 2 to 5 µm. This assures a compact and homogeneous
packing of the column which allows a regular circulation of the mobile phase without the
formation of preferential routes in the column (Figure 8.8).
It is prepared under conditions of controlled hydrolysis, by a sol–gel polymerization
of an alkoxysilicate (e.g. tetraethoxysilane) in the form of an emulsion, under the effect of
base-catalysed hydrolysis. Initially, tiny particles are formed 0.2 µm which grow in a
regular manner, by various methods, to form spheres that attain a few micrometres in
diameter.
Figure 8.8: Schematic structure of silica particle. [From R. E. Majors, LCGC May 1997, p. S8.]
8.3.2. Partition Chromatography (Liquid-Liquid chromatography)
The most widely used type of HPLC is partition chromatography, in which
the stationary phase is a second liquid that is immiscible with the liquid mobile
phase. In the past, most of the applications have been to nonionic, polar
compounds of low to moderate molecular mass (usually <3000).
The early forms of partition chromatography used liquid-liquid columns.
These have been replaced in modern LC systems by liquid bonded-phase columns.
In liquid-liquid chromatography, the liquid was held in place by physical
adsorption. In bonded-phase chromatography, on the other hand, it is attached by
chemical bonding, resulting in highly stable packings insoluble in the mobile
phase. Bonded-phase columns are also compatible with gradient elution
techniques.
Therefore, our discussion focuses exclusively on bonded-phase partition
chromatography.
Columns for Bonded-Phase Chromatography
The supports for the majority of bonded-phase packings for partition
chromatography are prepared from rigid silica, or silica-based, compositions.
These solids are formed as uniform, porous, mechanically sturdy particles
commonly having diameters of 1.5-10 µm with 3- and 5- µm particles being most
common. The surface of fully hydrolyzed silica (hydrolyzed by heating with 0.1 M
HCI for a day or two) is made up of chemically reactive silanol groups. Typical
silica surfaces contain about 81 µmol/m2 of OH groups and have surface areas of
100 to 300 m2/g.
The most useful bonded-phase coatings are siloxanes formed by reaction of
the hydrolyzed surface with an organochlorosilane. For example,
where R is an alkyl group or a substituted alkyl group.
Figure 8.9: The scheme of synthesizing bonded phase coatings

Normal- and Reversed-Phase Packings
Two types of partition chromatography are distinguishable based on the
relative polarities of the mobile and stationary phases. Early work in LC was based
on highly polar stationary phases such as triethylene glycol or water; a relatively
nonpolar solvent such as hexane or i-propyl ether then served as the mobile phase.
For historic reasons, this type of chromatography is now called normal-phase
chromatography.
In reversed-phase chromatography, the stationary phase is nonpolar, often a
hydrocarbon, and the mobile phase is a relatively polar solvent (such as water,
methanol, acetonitrile, or tetrahydrofuran).
In normal-phase chromatography, the least polar component is eluted first;
increasing the polarity of the mobile phase then decreases the elution time. With
reversed-phase chromatography, however, the most polar component elutes first,
and increasing the mobile-phase polarity increases the elution time.
Normal-phase partition chromatography and adsorption chromatography
overlap considerably. In fact, retention in most types of normal-phase
chromatography appears to be governed by adsorption-displacement processeses.
Bonded-phase packings are classified as reversed phase when the bonded
coating is nonpolar in character and as normal phase when the coating contains
polar functional groups. It has been estimated that more than three quarters of all
HPLC separations are currently performed in columns with reversed-phase
packings. The major advantage of reversed-phase separations is that water can be
used as the mobile phase. Water is an inexpensive, nontoxic, UV-transparent
solvent compatible with biological solutes. Also, mass transfer is rapid with
nonpolar stationary phases, as is solvent equilibration after gradient elution. Most
commonly, the R group of the siloxane in these coatings is a C8 chain (n-octyl) or
a C18chain (n-octyldecyl).
In most applications of reversed-phase chromatography, elution is carried out
with a highly polar mobile phase, for example, an aqueous solution containing
various concentrations of such solvents as methanol, acetonitrile, or
tetrahydrofuran. In this mode, care must be taken to avoid pH values greater titan
about 7.5 because the silica can form soluble silicate species, causing dissolution
of the stationary phase. Also, hydrolysis of the siloxane can occur in alkaline
solutions, which leads to degradation or dcstruction of the packing. Acid
hydrolysis of the siloxane can limit the pH to about 2.5 in acidic solutions.
In commercial normal-phase bonded packings, the R in the siloxane structure
is a polar functional group such as cyano, -C2H4CN; diol,
-C3H6OCH2CHOHCH2OH; amino, -C3H6NH2; and dimethylamino, C3H6N(CH3),.
The polarities of these packing materials vary over a considerable range, with the
cyano type being the least polar and the amino types the most. Diol packings are
intermediate in polarity. With normal-phase packings, elution is carried out with
relatively nonpolar solvents, such as ethyl ether, chloroform, and n-hexane.
8.3.3. lon-Pair Chromatography
Ion-pair chromatography, sometimes called paired-ion chromatography is a
subset of reversed-phase chromatography in which easily ionizable species are
separated on reversed-phase columns. In this type of chromatography. an organic
salt containing a large organic counter ion, such as a quaternary ammonium ion or
alkyl sulfonate, is added to the mobile phase as an ion-pairing reagent. Two
mechanisms for separation are postulated. In the first, the counter ion forms an
uncharged ion pair with a solute ion of opposite charge in the mobile phase. This
ion pair then partitions into the nonpolar stationary phase, giving differential
retention of solutes based on the affinity of the ion pair for the two phases.
Alternatively, the counter ion is retained strongly by the normally neutral
stationary phase and imparts a charge to this phase. Separation of organic solute
ions of the opposite charge then occurs by formation of reversible ion-pair
complexes with the more strongly retained solutes forming the strongest
complexes with the stationary phase. Some unique separations of both ionic and
nonionic compounds in the same sample can be accomplished by this form of
partition chromatography.
Table 8.1: Some example of paired –ion chromatography
pH Mobile phase Counter ion Sample

pH 9 Cyclohexane/CHCl3/1-pentanol N,N-dimethyl- carboxylic acids
protrityline
0.1MHClO4 Mixtures of tributyl phosphate,
ethyl acetate, butanol, ClO4- amines
CH2Cl2,and/or hexane
buffer HPO42- /PO43- Butanol/CH2Cl2/Hexane (Butyl)4N+ carboxylic acids
0.1MHClO4 CH2Cl2/ CHCl3/ butanol và /hoặc ClO4- amines
pH 5-6 Pentanol/CH2Cl2 và/hoặc CHCl3 Picrate amines
pH 6-8.5 Butanol/heptane (Butyl)4N+ sulfoamit
0.2-0.25M HClO4 Butanol/CH2Cl2/Hexane ClO4- Amines and

quaternary amines
0.1M methane-
sunfonic acid Butanol/CH2Cl2/Hexane CH3SO3- amines
pH 8.3 Butanol/CH2Cl2/Hexane (Butyl)4N+ carboxylic acids

8.3.4. Ion Chromatography
In ion-exchange chromatography, retention is based on the attraction between
solute ions and charged sites bound to the stationary phase. In anion exchangers,
positively charged groups on the stationary phase attract solute anions. Cation
exchangers contain covalently bound, negatively charged sites that attract solute
cations. Figure 28-22 shows the structure of a strong acid resin. Note the cross-
linking that holds the linear polystyrene molecules together. The other types of
resins have similar structures except for the active functional group.
Stationary phase - Ion Exchangers
Resins are amorphous (noncrystalline) particles of organic material.
Polystyrene resins for ion exchange are made by the copolymerization of styrene
and divinylbenzene (Figure 8.10). Divinylbenzene content is varied from 1 to 16%
to increase the extent of cross-linking of the insoluble hydrocarbon polymer. The
benzene rings can be modified to produce a cation-exchange resin, containing
sulfonate groups (–SO3-), or an anion-exchange resin, containing ammonium
groups (-NR3+) . If methacrylic acid is used in place of styrene, a polymer with
carboxyl groups results.
Figure 8.10: The structure of ion exchangers based on Polystyrene resins

Ion exchangers are classified as strongly or weakly acidic or basic. Sulfonate groups
of strongly acidic resins remain ionized even in strongly acidic solutions. Carboxyl
groups of the weakly acidic resins are protonated near pH 4 and lose their cation-
exchange capacity. “Strongly basic” quaternary ammonium groups (which are not really
basic at all) remain cationic at all values of pH. Weakly basic tertiary ammonium (-
CH2NHR2+) anion exchangers are deprotonated in moderately basic solution and lose
their ability to bind anions.
Ion-Exchange Selectivity
Consider the competition of Na+ and Li+ for sites on the cation-exchange
resin R-:
The equilibrium constant is called the selectivity coefficient, because it

describes the relative selectivity of the resin for Li + and Na+. Selectivities of
polystyrene resins tend to increase with the extent of cross-linking, because the
pore size of the resin shrinks as cross-linking increases. Ions such as Li +, with a
large hydrated radius, do not have as much access to the resin as smaller ions, such
as Cs+, do.
In general, ion exchangers favor the binding of ions of higher charge,
decreased hydrated radius, and increased polarizability. A fairly general order of
selectivity for cations is
Mobile phases
Depending upon the type of stationary phase, the counter ions present in the mobile
phase derived from acids (perchloric, benzoic, phthalic, methane sulfonic), or bases (the
most popular for anion analyses are variants of sodium hydroxide and sodium
carbonate/bicarbonate).
Figure 8.11: Elution of lanthanide(III) ions from a cation-exchange resin. Eluent was varied from
20 to over 25 min.
Conductivity detectors
Besides the spectrophotometric detectors based on absorbance or fluorescence of
uv/visible radiation, and used when the mobile phase does not absorb appreciably,
another mode of detection exists based upon electrolyte conductivity.
Thus, at the outlet of the column, the conductance (the inverse of the
resistance) of the mobile phase is measured between two microelectrodes. The
measuring cell should be of a very small volume (approx. 2µL).
The difficulty is to recognize in the total signal the part due to ions or ionic
substances present in the sample. In order to do direct measurements, the ionic
charge of the mobile phase has to be as low as possible and the measuring cell
requires strict temperature control to within 0.01 oC because of the high
dependence of conductance on temperature (∼5%/).
Ion suppressors
The mobile phase contains ions that create a background conductivity, making it
difficult to measure the conductivity due only to the analyte ions as they exit the column.
To improve the signal to noise ratio, when using a conductivity detector, a device called a
suppressor, designed to selectively remove the mobile phase ions is placed after the
analytical column and before the detector. The principle consists to convert the mobile
phase ions to a neutral form or replacing them by others of higher conductivity (Figure
8.12).
.
Figure 8.12: Schematic illustrations of (a) suppressed-ion anion chromatography and (b)
suppressed-ion cation chromatography.
8.3.5. Molecular Exclusion Chromatography
In molecular exclusion chromatography (also called size exclusion or gel filtration
or gel permeation chromatography), molecules are separated according to size. Small
molecules penetrate the pores in the stationary phase, but large molecules do not. Because
small molecules must pass through an effectively larger volume, large molecules are
eluted first. This technique is widely used in biochemistry to purify macromolecules.
Size-exclusion. or gel, chromatography is a powerful technique that is particularly
applicable to highmolecular-mass species." Packings for size-exclusion chromatography
consist of small (-10 flm) silica or polymer particles containing a network of uniform
pores into which solute and solvent molecules can diffuse. While in the pores, molecules
are effectively trapped and removed from the flow of the mobile phase. The average
residence time in the pores depends on the effective size of the analyte molecules.
Molecules larger than the average pore size of the packing are excluded and thus suffer
essentially no retention; such species are the first to be eluted. Molecules having
diameters significantly smaller than the pores can penetrate or permeate throughout the
pore maze and are thus entrapped for the greatest time; these are last to be eluted.
Between these two extremes are intermediate size molecules whose average penetration
into the pores of the packing depends on their diameters.
Within this group, fractionation occurs, which is directly related to molecular size
and to some extent molecular shape. Note that size-exclusion separations differ from the
other procedures we have been considering in that no chemical or physical interaction
between an-alytes and the stationary phase are involved.
Column packing
Two types of packing for size-exclusion chromatography are encountered: polymer
bcads and silica-based particles, both of which have diameters of 5 to 10 µm.
Most early size-exclusion chromatography was carried out; n cross-linked styrene-
divinylbenzene copolymers similar in structure (except that the sulfonic acid groups arc
absent) to that shown in Figure 8.10. The pore size of these polymers is controlled by the
extent of cross-linking and hence the relative amount of divinylbenzene present during
manufacture. Polymeric packings having several different average pore sizes are thus
marketed. Originally, styrene-divinylbenzene gels were hydrophobic and thus could be
used only with nonaqueous mobile phases. Now, however, hydrophilic gels are available,
making possible the use of aqueous solvents for the separation of large, water-soluble
molecules such as sugars. These hydrophilic gels are sulfonated divinylbenzenes or
polyacrylamides (Figure 8.14). Chromatography based on the hydrophilic packings was
once called gel filtration, and techniques based on hydrophobic packings were termed gel
permeation. Today, both techniques are described as size-exclusion methods.
Figure 8.13: Applications of size –exclusion chromatography

Figure 8.14: Structure of polyacrylamide
8.4. Thin layer chromatography (planar chromatography)
8.4.1. Principle of TLC
The principle of the separation between phases of the sample components is similar
to that of HPLC, though the migration of the constituents through the stationary phase is
different. Separation is conducted on a thin layer (100–200 µm) of stationary phase,
usually based upon silica gel and deposited on a rectangular plate made out of glass,
plastic or aluminium of a few centimetres in dimensions.
To maintain the stationary phase on the support and to assure the cohesion of the
particles, an inert binder like gypsum (or organic linker) is mixed into the stationary
phase during the manufacture of the plate. The constituents can be identified by
simultaneously running standards with the unknown.
There are three steps to conduct a separation with this technique:
Deposition of the sample
A small volume of sample (between a few nanolitres to a few microlitres), dissolved
in a volatile solvent, is deposited close to the bottom of the plate as a small spot of
about 1–2mm in diameter. This deposit is made either manually, or automatically, with a
flat ended capillary (Figure 8.15). The spot can also have the form of a horizontal band of
a few millimetres which is obtained by automatic spraying of the sample. This last
method has the advantage of having a high reproducibility, indispensable of course, for
quantitative analysis. The prepared plate is then placed in a glass developing chamber
that contains a small amount of the appropriate developing solvent. The chamber is then
covered (Figure 8.15). The position at which the sample has been deposited must be
above the level of the solvent.
Developing the plate
The mobile phase rises up the stationary phase by capillarity, moving the
components of the sample at various rates because of their different degrees of interaction
with the matrix and solubility in the solvent. Their separation may be complete in a few
minutes. When the solvent front has travelled a sufficient distance (several centimetres),
the plate is withdrawn from the chamber, the position attained by the mobile phase is
immediately noted then it is evaporated.
Figure 8.15: Developing chamber and TLC plate. Left, available in a variety of dimensions
According to the size of the plates (of 5×5to20×20 cm), the chambers are made of
glass and equipped with a tight fitting cover. Right, typical appearance after the TLC
plate is partially dry. A faint line can be observed at the location of the solvent upper
limit. This line is called the solvent front. Calculation of R f (cf. paragraph 5.4). A
substance that does not migrate from the sample origin has a R =0.
8.4.2. Identifying the spots

The localization of each component on the plate, which has now lost all of the
eluent, consists to measure their migration distance from the original deposition. To
locate colourless compounds the plate must be developed (Figure 5.2). In order to
facilitate the visualization of the spots,manufacturers sell plates that contain a fluorescent
salt of zinc which emits a bright green fluorescence when the plate is irradiated with a
UV mercury vapour lamp = 254 nm and observed in a viewing cabinet. All compounds
absorbing at this wavelength appear as a dark spot (or sometimes coloured) against a
bright green blue background. Another method to make compounds visible, almost
universally used, consists of heating the plate after spraying itwith sulfuric acidwhich
leaves charred blots behind. This approach is, however, not adapted to quantitative TLC;
in this case, development is effected by immersion of the plate, using a general
(phosphomolybdic acid, vanillin), or specific (e.g. ninhydrin in alcoholic solution for
amino acids) reagent. Hundreds of reagents have been described that serve to introduce
chromophores or fluorophores groups into the analytesmolecules after separation.
The use of square TLC plates allows two-dimensional chromatography to be carried
out using two successive elutions with two different mobile phases (Figure 5.3). A high
degree of separation can be achieved.
8.4.3. Quantitative TLC
In order to use TLC as a quantitative method of analysis, it is essential to
quantify the spots (Figures 5.1 and 5.6), along with definitions for all of the usual
parameters (specificity, range of the domain of linearity, precision, etc.). This is
done by placing the plate under the lens of a densitometer (or scanner) that can
measure either absorption or fluorescence at one or several wavelengths. This
instrument produces a pseudo-chromatogram that contains peaks whose areas can
be measured. In fact it is actually an isochronic image of the separation at the final
instant. In TLC a spot is usually detectable if it corresponds at least to a few ng of
a compound UV absorbent.
Problems
8-1) When considering the van Deemter equation, why does HPLC require small column
packing particles?
8-2) What is the order in which acetone, acetamide and 1,2-dichloroethane would be
eluted from an HPLC column packed with C18 hydrocarbon bonded to the siloxane
backbone material?
8-3) What is a guard column and why is it used in HPLC but not in GC systems?
8-4) Consider the following three compounds HOCH 2CH2OH(ethylene glycol),
CH3CH2CH2CH2CH2CH3(hexane), C6H5CH3 (toluene), where toluene has a polarity
that is in between the two other compounds.
a) In what order would the compounds elute in normal-phase HPLC? For a given
mobile phase composition, the retention time of ethylene glycol is 5.6 minutes. Will
the retention time increase or decrease by increasing the polarity of the mobile
phase?
b) In what order would the compounds elute in reversed-phase HPLC? For a given
mobile phase composition, the retention time of hexane is 5.6 minutes. Will the
retention time increase or decrease by increasing the polarity of the mobile phase?
8-5) Predict the elution order of the following solutes in reversed phase HPLC.
8-6) An HPLC analysis was conducted for caffeine on a sports drink. A 10.1 ppm
methanol IS standard was introduced both into the sample and a standardized
solution of 304 ppm of caffeine. The measured by a diode‐array detector at each
λmax for the absorptions for methanol and for caffeine are summarized in the table
below. What is the concentration of caffeine in that sports drink?
IS Caffeine
Sample 23141 52777
304 ppm Caffeine standard 28441 77313
8-7) The concentration of PAHs in soil can be determined by first extracting the PAHs
with methylene chloride. The extract is then diluted, if necessary, and the PAHs are
separated by HPLC using a UV/Vis or fluorescence detector. Calibration is achieved
using one or more external standards. In a typical analysis, a 2.013-g sample of
dried soil is extracted with 20.00 mL of methylene chloride. After filtering to
remove the soil, a 1-mL portion of the extract is removed and diluted to 10 mL with
acetonitrile. Injecting 5 µL of the diluted extract into an HPLC gives a signal of
0.217 (arbitrary units) for the PAH fluoranthene. When 5 µL of a 20.0-ppm
fluoranthene standard is analyzed using the same conditions, a signal of 0.258 is
measured. Report the parts per million of fluoranthene in the soil.
8.8) The amount of caffeine in an analgesic tablet was determined by HPLC using a
normal calibration curve. Standard solutions of caffeine were prepared and analyzed
using a 10-µL fixed-volume injection loop. Results for the standards are
summarized in the following table.
The sample was prepared by placing a single analgesic tablet in a small beaker and
adding 10 mL of methanol. After allowing the sample to dissolve, the contents of
the beaker, including the insoluble binder, were quantitatively transferred to a 25-
mL volumetric flask and diluted to volume with methanol. The sample was then
filtered, and a 1.00-mL aliquot was transferred to a 10-mL volumetric flask and
diluted to volume with methanol. When analyzed by HPLC, the signal for the
caffeine was found to be 21469. Report the number of milligrams of caffeine in the
analgesic tablet.
8-9) The concentrations of Cl–, NO3–, and SO42– may be determined by ion
chromatography. A 50-µL standard sample of 10.0-ppm Cl –, 2.00-ppm NO3–, and
5.00-ppm SO42– gave signals (in arbitrary units) of 59.3, 16.1, and 6.08, respectively.
A sample of effluent from a wastewater treatment plant was diluted tenfold, and a
50-µL portion gave signals of 44.2 for Cl–, 2.73 for NO3–, and 5.04 for SO42–. Report
the parts per million of each anion in the effluent sample.
8-10) A series of polyvinylpyridine standards of different molecular weight were
analyzed by size-exclusion chromatography, yielding the following results.
When a preparation of polyvinylpyridine of unknown formula weight was analyzed

the retention volume was found to be 8.45. Report the average formula weight for
the preparation.
TÀI LIỆU THAM KHẢO

1. Modern analytical chemistry, David Harvey, Copyright © 2000 by The
McGraw-Hill Companies, Inc.
2. Analytical Electrochemistry, Joseph Wang, 2000, Second Edi., Wiley
3. Chromatographic methods, A.Braithwaithe and F.J. Smith, 1999, Kluwer
Academic Publishers.
4. Modern practice of Gas chromatography, Robert L. Grob, 1995, Wiley-
Interscience Publication.
5. Practical HPLC methodology and applications,Brian A. Bidlingmeyer,
1992, Wiley-Interscience Publication.
6.. Quantitative chemical analysis, Deniel C. Harris, 1999, W.H Freeman and
company

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Bùi Xuân Vững Instrumental Analysis Chemistry

PART A: SPECTROCHEMICAL METHODS OF ANALYSIS

magnetic component of electromagnetic radiation is responsible for the absorption of

ν = C/λ = 3.00 x 108 m.s-1/(325 x 10-2 m) = 9.23 x 107 s-1

Wavenumber is widely used in infrared vibrational spectroscopy. Wavenumber is

Diffraction, 0.1-100 Ao - Nonbonding electrons

ultraviolet and visible

Chapter 1: ULTRAVIOLET-VISIBLE (UV-VIS) ABSORPTION

1.1. UV-VIS Absorption of molecules

 Energy of hν of a photon is equal to the difference between the two molecular

1.2. Beer’s law

According to Beer’s law: A = εlC

c. Effects of polychromatic radiation

Figure 1.3: The effect of polychromatic

d. Effect of the stray radiation

The radiation used for absorbance measurements is usually contaminated

Based on such measuring principle, a UV-VIS spectrometer needs to have the

Figure 1.4: Some types of UV-VIS cuvettes

is required for constan intensities. This lamp provides a useful continuous

Figure 1.5: Mechanism of diffraction from an echellette-type grating

e. Designs of UV-VIS spectrometers

Single beam instrument:

Figure 1.8: a) a typical single beam spectrometer, b) its optical diagram

Figure 1.8bis: Structural diagram of a double beam spectrometer.

1.4. Quantitative analysis applications

 Good accuracy: The relative error in concentration for a analytical procedure

 Use a suitable pipette to add 10.0 mL of a water sample to to each of 5

1.4.3. Analysis of Mixtures

individual components in a mixture even if there is strong overlap in their

curve, the concentration of the solution of unknown X concentration. Construct

Measurements were made in 1.0 cm glass cells.

Chapter 2: ATOMIC EMISSION SPECTROSCOPY (AES)

Figure 2.1: A simple scheme of electron layers and their energies

the arrangement of the shells from

Figure 2.2: Energy scheme of subshells

These excited atoms are in 2.3). So an emission spectrum is one

The more outmost electrons the

Figure 2.3: Emission spectrum of Na

According to the Einstein equation, we have:

intensity Il and concentration C is illustrated by a plot in Figure 2.4.

where T is temperature (K) and k is Boltzmann’s constant (1.381x10-23 J/K).

measuring rate, or store, print, compare results, an atomic emission spectrometer is

Oxidant Fuel Temperature

The flame temperature depends upon:

Flame consists of three regions as

Figure 2.6: Diagram of a flame spectrometer using to determine alkali metals

Figure 2.7: Inductively coupled plasma source

In reality ICP is a thermoelectric energy source of a flame maintained by an

In the method of using ultrasonic nebulizer, sample solution is directed onto a

Figure 2.9: A conventional nebulizer with carrier argon

Figure 2.10: An ultrasonic nebulizer

Figure 2.11: Overview of a Basic Inductively Coupled Plasma—Atomic Emission

2.6. Data Collection

2.7.2. Standard Addition Method

advantage of AES is the linearity of response over a very large range of

2-2) To measure Ca in breakfast cereal, 0.5216 g of crushed Cheerios was ashed

Determine the parts per million Na in the sample of oat bran.

What is the %w/w Pb in a sample of brass that gives an emission intensity of

Chapter 3: ATOMIC ABSORPTION SPECTROSCOPY (AAS)

ΔE = (Em – Eo) = hγ (3.1)

absorbed, we have an absorption line with the wavelength of li , meaning that

Figure 3.1: AAS and AES processes of Na atoms

remain to depend upon concentration N of the atom in the absporption medium.